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# 1 Mass vs. angular momentum and parity (𝐽^𝑃) for neutral kaon resonances. Arrows indicate resonances accessible by Primakoff excitation.

RUTGERS-01-88

Radiative Decay Width Measurements of Neutral Kaon Excitations Using the Primakoff Effect

A. Alavi-Harati12, T. Alexopoulos12, M. Arenton11, K. Arisaka2, S. Averitte10, R.F. Barbosa7,∗∗, A.R. Barker5, M. Barrio4, L. Bellantoni7, A. Bellavance9, J. Belz10, D.R. Bergman10, E. Blucher4, G.J. Bock7, C. Bown4, S. Bright4, E. Cheu1, S. Childress7, R. Coleman7, M.D. Corcoran9, G. Corti11, B. Cox11, A. Cunha10, A.R. Erwin12, R. Ford7, A. Glazov4, A. Golossanov11, G. Graham4, J. Graham4, E. Halkiadakis10, J. Hamm1, K. Hanagaki8, S. Hidaka8, Y.B. Hsiung7, V. Jejer11, D.A. Jensen7, R. Kessler4, H.G.E. Kobrak3, J. LaDue5, A. Lath10, A. Ledovskoy11, P.L. McBride7, D. Medvigy10, P. Mikelsons5, E. Monnier4,∗, T. Nakaya7, K.S. Nelson11, H. Nguyen7, V. O’Dell7, R. Pordes7, V. Prasad4, X.R. Qi7, B. Quinn4, E.J. Ramberg7, R.E. Ray7, A. Roodman4, S. Schnetzer10, K. Senyo8, P. Shanahan7, P.S. Shawhan4, J. Shields11, W. Slater2, N. Solomey4, S.V. Somalwar10,†, R.L. Stone10, E.C. Swallow4,6, S.A. Taegar1, R.J. Tesarek10, G.B. Thomson10, P.A. Toale5, A. Tripathi2, R. Tschirhart7, S.E. Turner2, Y.W. Wah4, J. Wang1, H.B. White7, J. Whitmore7, B. Winstein4, R. Winston4, T. Yamanaka8, E.D. Zimmerman4

1 University of Arizona, Tucson, Arizona 85721

2 University of California at Los Angeles, Los Angeles, California 90095

3 University of California at San Diego, La Jolla, California 92093

4 The Enrico Fermi Institute, The University of Chicago, Chicago, Illinois 60637

5 University of Colorado, Boulder, Colorado 80309

6 Elmhurst College, Elmhurst, Illinois 60126

7 Fermi National Accelerator Laboratory, Batavia, Illinois 60510

8 Osaka University, Toyonaka, Osaka 560-0043 Japan

9 Rice University, Houston, Texas 77005

10 Rutgers University, Piscataway, New Jersey 08854

11 The Department of Physics and Institute of Nuclear and Particle Physics, University of Virginia, Charlottesville, Virginia 22901

12 University of Wisconsin, Madison, Wisconsin 53706

Permanent address C.P.P. Marseille/C.N.R.S., France

∗∗Permanent address University of São Paulo, São Paulo, Brazil

To whom correspondence should be addressed. Electronic address: somalwar@physics.rutgers.edu

The KTeV Collaboration

## Abstract

We produce a sample consisting of 147 candidate events, with minimal backgrounds, of the mixed axial vector pair ($`K_1(1270)`$-$`K_1(1400)`$) by exciting $`K_L`$’s in the Coulomb field of lead and report the first measurements of the radiative widths $`\mathrm{\Gamma }_r(K_1(1400))=280.8\pm 23.2(stat)\pm 40.4(syst)`$ keV and $`\mathrm{\Gamma }_r(K_1(1270))=73.2\pm 6.1(stat)\pm 28.3(syst)`$ keV. We also place 90% CL upper limits $`\mathrm{\Gamma }_r(K^{}(1410))52.9`$ keV for the vector state and $`\mathrm{\Gamma }_r(K_2^{}(1430))5.4`$ keV for the tensor state. These measurements allow for significant tests of quark-model predictions of radiative widths for the low-lying vector mesons.

Several resonant excitations of the neutral kaon are known to exist , most having been observed indirectly using partial wave analysis . Figure 1 is a schematic representation of the neutral kaon excitations with central masses less than 1.5 GeV/c2. The axial vector pair ($`K_1(1270)`$-$`K_1(1400)`$) is interesting because it is a (coherent) mixture of the singlet $`{}_{}{}^{1}P_{1}^{}`$ and the triplet $`{}_{}{}^{3}P_{1}^{}`$ states , parameterized by the mixing angle $`\mathrm{\Theta }`$: $`K_1(1270)={}_{}{}^{3}P_{1}^{}\mathrm{sin}\mathrm{\Theta }+{}_{}{}^{1}P_{1}^{}\mathrm{cos}\mathrm{\Theta }`$ and $`K_1(1400)={}_{}{}^{3}P_{1}^{}\mathrm{cos}\mathrm{\Theta }+{}_{}{}^{1}P_{1}^{}\mathrm{sin}\mathrm{\Theta }`$. The radiative decay widths of the kaon excitations, $`\mathrm{\Gamma }_r(K^{})=\mathrm{\Gamma }(K^{}K+\gamma )`$, are sensitive to the magnetic moments of the constituent quarks . Radiative widths have been calculated for low-lying mesons using both a dynamic quark model  and a relativistic quark model . Experimentally, only $`\mathrm{\Gamma }_r(K^{}(892))`$ has been measured  so far.

The Primakoff effect , i.e. excitation by the Coulomb field, can be used to measure radiative widths since it is the inverse of radiative decay. In this Letter, we use the full dataset collected during the 1996-97 run of the KTeV experiment at Fermilab to study Primakoff production in two channels: the six-body $`K^{}(892)`$$`\pi ^0`$ channel, exemplified by $`K^{}(1410)`$ or $`K_1(1400)`$ $`K^{}(892)\pi ^0[K_S\pi ^0]\pi ^0[(\pi ^+\pi ^{})(\gamma \gamma )](\gamma \gamma )`$, which has two $`\pi ^0`$’s, and the four-body $`K_S`$$`\pi ^0`$ channel, exemplified by $`K^{}(892)`$ or $`K^{}(1410)`$ $`K_S\pi ^0(\pi ^+\pi ^{})(\gamma \gamma )`$, which has a single $`\pi ^0`$. In the $`K^{}(892)`$$`\pi ^0`$ channel, we observe 147 candidate events which are predominantly the axial vector $`K_1(1400)`$ with a small admixture of $`K_1(1270)`$. Using a large sample of $`K^{}(892)`$’s from the $`K_S`$$`\pi ^0`$ channel for normalization, we report the first measurements of the radiative widths for the axial vector pair. We also use the $`K_S`$$`\pi ^0`$ channel to place the first upper limit on $`\mathrm{\Gamma }_r(K^{}(1410))`$ and a stringent upper limit on $`\mathrm{\Gamma }_r(K_2^{}(1430))`$.

For high particle energies and small production angles, the rate of exciting a $`K_L`$ to a $`K^{}`$ in the Coulomb field of a nucleus A is given by

$$\frac{d\sigma }{dt}(K_L+AK^{}+A)=\pi \alpha Z^2\left(\frac{2S_K^{}+1}{2S_K+1}\right)\frac{\mathrm{\Gamma }(K^{}K+\gamma )}{k^3}\frac{t^{}}{t^2}|f_{EM}|^2,$$
(1)
where $`\alpha `$ is the fine structure constant, $`Z`$ is the atomic number of the nucleus, $`S_K`$ and $`S_K^{}`$ are the spins of $`K_L`$ and the resonance, respectively, $`k=(m_K^{}^2m_K^2)/2m_K^{}`$, $`t`$ is the magnitude of the square of the momentum transfer and $`t^{}=tt_{min}`$, $`\sqrt{t_{min}}=(m_K^{}^2m_K^2)/2P_K`$, where $`P_K`$ is the laboratory momentum of the $`K_L`$. Finally, $`f_{EM}`$ is the nuclear electric form factor. Thus, the rate of Primakoff production is directly proportional to the radiative width.

KTeV utilized an $`800`$ GeV/c proton beam to generate two neutral beams consisting kaons, neutrons and some hyperons. In the E832 configuration , one of the beams passed through a regenerator which was located $``$124 m from the target. The regenerator consisted of 84 modules of 2cm-thick plastic scintillator followed by a module composed of a lead-scintillator sandwich. Since the Primakoff effect is proportional to $`Z^2`$ of the target material, more than 98%  of the observed Primakoff excitations (equation 1) were produced in the final lead pieces. The regenerator was instrumented with photomultiplier tubes which enabled us to tag and reject backgrounds from inelastic interactions. We detect $`\pi ^+`$$`\pi ^{}`$ tracks from $`K_S`$ decays using a drift chamber spectrometer system and photons from $`\pi ^0`$ decays using a pure CsI electromagnetic calorimeter. The event trigger was initiated by signals from two scintillator hodoscopes located downstream of the spectrometer and required hits in the drift chambers consistent with two oppositely charged tracks. The decay volume was surrounded by a near-hermetic set of devices to veto photons.

In the offline analysis, the fiducial region for the decay vertex of $`K_S\pi ^+\pi ^{}`$ is restricted to 15 m downstream of the regenerator. We reconstruct $`\pi ^0`$’s using pairs of energy clusters in the calorimeter. The clusters are required to have energies greater than 1 GeV and photon-like spatial distributions. To reject electrons, we require that the ratio of energy deposited in the calorimeter to the particle momentum as measured by the spectrometer be $`<0.8`$.

To reconstruct the $`K^{}(892)K_S\pi ^0`$ decays used for normalization, in the four-body channel we require the invariant masses of the $`\gamma \gamma `$ and the $`\pi ^+`$$`\pi ^{}`$ to be within 10 MeV/c2 of the $`\pi ^0`$ and $`K_S`$ invariant masses, respectively. We isolate Primakoff (forward) production by demanding that the square of the transverse momentum ($`p_t^2`$) of the $`\pi ^+\pi ^{}\gamma \gamma `$ with respect to a line connecting the target and the decay vertex of $`K^{}(892)`$ be less than $`0.001(GeV/c)^2`$. We further require $`\pi ^+`$$`\pi ^{}`$ $`p_t^2`$ $`>0.01(GeV/c)^2`$ because the daughter $`K_S`$ recoils against the $`\pi ^0`$. The resulting sample of 29,399 $`K^{}(892)K_S\pi ^0`$ decays with $`K_S`$ energy between 30 and 210 GeV, and the $`K^{}(892)`$ energy between 55 and 225 GeV is shown in figure 2(top).

The requirements for the $`K^{}(892)`$$`\pi ^0`$ six-body channel are similar, except for changes to account for the extra $`\pi ^0`$ and differences in kinematics. The photon pairings for the two $`\pi ^0`$’s are determined using a $`\chi ^2`$ formed by comparing $`M_{\gamma \gamma }`$ and $`M_{\pi ^+\pi ^{}\gamma \gamma }`$ to the known masses of $`\pi ^0`$ and $`K^{}(892)`$, respectively. The $`K_S`$$`\pi ^0`$ mass for the daughter $`K^{}(892)`$ is required to be within 101 MeV/c2 (two mass-widths) of the $`K^{}(892)`$ mass and its $`p_t^2`$ to be $`>0.03(GeV/c)^2`$. The $`p_t^2`$ cut also serves to eliminate background from Primakoff-produced $`K^{}(892)`$’s when accompanied by an accidental $`\pi ^0`$. To eliminate events in which two kaons decay to a charged and a neutral pion pair, we remove events for which the four-photon invariant mass is within 20 MeV of the $`K_L`$ mass. The resulting sample of ($`K_S`$$`\pi ^0`$)$`\pi ^0`$ events with total energy greater than 90 GeV is depicted in figure 3 and shows a clustering near 1.4 GeV/c2. The mass projection shows the resonant signature exhibited by events with $`p_t^2<0.001(GeV/c)^2`$ and the $`p_t^2`$ projection shows the sharp fall-off confirming Primakoff production. There are 147 events within the mass fiducial region (1.1-1.64 GeV/c2). Figure 4 shows the invariant mass and $`p_t^2`$ of the daughter $`K^{}(892)`$ where the $`p_t^2`$ displays a Jacobian distribution expected of a daughter particle in a two-body decay.

The possible candidates for the observed $`K^{}(892)`$$`\pi ^0`$ resonance are $`K_0^{}(1430)`$, K(1460), $`K^{}(1410)`$, $`K_2^{}(1430)`$, $`K_1(1270)`$, and $`K_1(1400)`$ (figure 1). The selectivity of the Primakoff effect rules out $`K_0^{}(1430)`$ and K(1460) because of spin-parity conservation and the J=0$`\to ̸`$ J=0 selection rule, respectively. Contributions from the vector $`K^{}(1410)`$ and tensor $`K_2^{}(1430)`$ can be eliminated because both have significant branching fractions to $`K_S`$$`\pi ^0`$ , yet we see no evidence for their presence in this ($`K_S`$$`\pi ^0`$) channel; note the lack of resonance near 1.4 GeV/c2 in figure 2 (bottom). We fit a combination of $`K^{}(892)`$ and $`K^{}(1410)`$($`K_2^{}(1430)`$) simulations to the data and confirm that the signal from $`K^{}(1410)`$($`K_2^{}(1430)`$) in the $`K_S`$$`\pi ^0`$ channel is consistent with zero: $`4.0\pm 6.0(0.1\pm 3.8)`$ $`K^{}(1410)`$($`K_2^{}(1430)`$) events. Using the known branching fractions  of $`K^{}(1410)`$ and $`K_2^{}(1430)`$ to $`K_S`$$`\pi ^0`$ and $`K^{}(892)`$$`\pi ^0`$, we translate these results into a negligible $`2.4\pm 3.6(0.0\pm 0.7)`$ event contribution of $`K^{}(1410)`$($`K_2^{}(1430)`$) in the $`K^{}(892)`$$`\pi ^0`$ channel. Thus we are left with only the axial vector pair ($`K_1(1270)`$-$`K_1(1400)`$) as a possible candidate for the observed resonance in the $`K^{}(892)`$$`\pi ^0`$ channel.

We cross-checked the axial vector nature of the observerd signal using the distributions of Gottfried-Jackson (GJ) angles222In the GJ frame the excited kaon is at rest at the origin, the z-axis is along the momentum of the incoming $`K_L`$ and the y-axis is perpendicular to the production plane . $`\theta `$ and $`\varphi `$. These distributions generally confirm our axial vector assignment. However, due to the relatively strong angular dependence of the detector acceptance, they do not have strong discrimination power between the axial vector pair and $`K^{}(1410)`$ and $`K_2^{}(1430)`$.

We now compute the radiative widths for the axial vector pair. It is difficult to decompose the mass spectrum of the observed signal into $`K_1(1270)`$ and $`K_1(1400)`$ because their mass separation is comparable to their widths. Nonetheless, mass information alone tells us that the contribution from $`K_1(1270)`$ is slight: only $`8.8\pm 8.6`$ events are due to $`K_1(1270)`$. However, a significantly better resolution is possible because the Primakoff effect can produce only the singlet ($`{}_{}{}^{1}P_{1}^{}`$) component of the axial vector pair  and the singlet-triplet mixing angle $`\mathrm{\Theta }`$ has been measured. Using $`\mathrm{\Theta }=56\pm 3^{}`$  together with the known branching ratios of $`K_1(1270)`$ and $`K_1(1400)`$, we resolve the observed signal into $`11.4\pm 1.0(stat)\pm 4.1(extsyst)`$ $`K_1(1270)`$ events and $`134.4\pm 11.1(stat)4.1(extsyst)`$ $`K_1(1400)`$ events, where the (external) systematic error is due to the measurement uncertainties in the mixing parameter $`\mathrm{\Theta }`$ and in the $`K_1(1270)`$ and $`K_1(1400)`$ branching fractions to the $`K^{}(892)`$$`\pi ^0`$ channel. This decomposition, depicted in figure 3, leads to $`\mathrm{\Gamma }_r(K_1(1270))=73.2\pm 6.1(stat)\pm 8.2(intsyst)\pm 27.0(extsyst)`$ keV and $`\mathrm{\Gamma }_r(K_1(1400))=280.8\pm 23.2(stat)\pm 31.4(intsyst)\pm 25.4(extsyst)`$ keV, where we have used our $`K^{}(892)`$ sample (figure 2) for normalization purposes since $`\mathrm{\Gamma }_r(K^{}(892))`$ is known experimentally . Our measurements share internal systematic errors of 8.7% due to uncertainties in the strong production (discussed below), 6.6% due to detector acceptance effects, and 2.4% due to the 3.6 event uncertainty in the possible contributions from $`K^{}(1410)`$ and $`K_2^{}(1430)`$, as discussed earlier. The uncertainty in the $`K^{}(892)`$ radiative width measurement  causes an additional 8.5% (external) systematic error.

Primakoff production is characterized by a sharp ($`t^1`$) forward production (equation 1) allowing a strict $`p_t^2<0.001(GeV/c)^2`$ cut which virtually eliminates all potential backgrounds; see figure 3. Based on an extrapolation from the large $`p_t^2`$ ($`>0.1(GeV/c)^2`$) region, we estimate 1.2 events out of 147 signal candidate events to be due to incoherent production and other possible backgrounds such as those from the decay products of the $`\mathrm{\Lambda }`$’s and $`\varphi `$’s produced when neutrons in the beam interact with the regenerator.

Coherent strong production and its interference (with unknown strength) with Primakoff production are expected to be small at our energies . Indeed, a maximum likelihood fit in the $`p_t^2`$ variable for the strong production and the strength of the strong-Coulomb interference using the prescription given in  indicates that the strength of interference preferred by our data is consistent with zero. A constructive (destructive) interference would mean that the actual number of Primakoff events is less (more) than what we observe. The mean change in our estimate of Primakoff production corresponding to one standard deviation variation in the interference strength is 8.7%, which we have taken to be the systematic error due to the uncertainties in strong production.

Earlier, we used the absence of a resonance in the $`K_S`$$`\pi ^0`$ channel at $``$1.4 GeV/c2(figure 2) to limit the $`K^{}(1410)`$ and $`K_2^{}(1430)`$ contributions to the observed ($`K_1(1270)`$-$`K_1(1400)`$) axial vector pair signal. A further benefit of this finding is that we are able to limit the radiative widths $`\mathrm{\Gamma }_r(K^{}(1410))`$ and $`\mathrm{\Gamma }_r(K_2^{}(1430))`$ to 52.9 and 5.4 keV, respectively, at 90% CL. $`\mathrm{\Gamma }_r(K^{}(1410))`$ has not been examined experimentally before, whereas $`\mathrm{\Gamma }_r(K_2^{}(1430))`$ was previously limited to 84 keV at 90% CL . The $`\mathrm{\Gamma }_r(K_2^{}(1430))`$ limit is far stricter than the $`\mathrm{\Gamma }_r(K^{}(1410))`$ limit principally because the branching fraction for $`K_2^{}(1430)K_S\pi ^0`$ is substantially larger  than the same for $`K^{}(1410)K_S\pi ^0`$.

The predicted radiative widths for the axial vector mesons , are 538 keV for $`K_1(1400)`$ and 175 keV for $`K_1(1270)`$; compare to our results, $`280.8\pm 46.6`$ keV and $`73.2\pm 28.9`$ keV, respectively. We note that the theoretical model in  is very sensitive to the quark masses ($`m_{u,s}`$) and rms momenta ($`\beta _{uu,us,ss}`$) of the quarks within the mesons. The predictions are based on certain choices for $`m`$ and $`\beta `$, but other choices with up to 30% variation in $`m`$ are possible.

Our 90% CL upper limit on the vector $`K^{}(1410)`$ radiative width is 52.9 keV. In the naive quark model, this state is the first radial excitation of $`K^{}(892)`$ and its radiative width calculation should be similar to that for $`K^{}(892)`$, for which $`\mathrm{\Gamma }_r(K^{}(892))=116.5\pm 9.9`$ keV . The smaller value for $`K^{}(1410)`$ may be due to a reduced overlap of the quark wavefunction for this higher radial excitation, but further guidance from theory is needed.

Finally, we have substantially improved the upper limit on the radiative width of the tensor $`K_2^{}(1430)`$ from 84 keV  to 5.4 keV (at 90% CL). Babcock and Rosner  used $`SU(3)`$ invariance to predict that excitations with $`J^{PC}=1^{++}`$ or $`2^{++}`$ would have vanishing radiative widths. In the limit of $`SU(3)`$, $`K_2^{}(1430)`$ has $`C=+1`$; thus, our limit lends support to Babcock and Rosner’s prediction  and serves as a direct test of the naive quark model and $`SU(3)`$-breaking.

We thank J. Bronzan for theoretical guidance and gratefully acknowledge the support and effort of the Fermilab staff and the technical staffs of the participating institutions for their vital contributions. This work was supported in part by the U.S. Department of Energy, The National Science Foundation and The Ministry of Education and Science of Japan. In addition, A.R.B., E.B. and S.V.S. acknowledge support from the NYI program of the NSF; A.R.B. and E.B. from the Alfred P. Sloan Foundation; E.B. from the OJI program of the DOE; K.H., T.N. and M.S. from the Japan Society for the Promotion of Science; and R.F.B. from the Fundação de Amparo à Pesquisa do Estado de São Paulo. P.S.S. acknowledges receipt of a Grainger Fellowship.


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# Partial observables

## I Introduction

The notion of “observable quantity”, or “observable”, plays a central role in many areas of physics. Roughly, observable quantities are the quantities involved in physical measurements. They give us information on the state of a physical system and may be predicted by the theory. In quantum mechanics, observables are represented by self-adjoint operators. In gauge theory, we make the distinction between gauge-invariant quantities, which correspond to observables, and gauge-dependent quantities, which do not not.

The notion of observable, however, raises a certain number of issues, which have generated discussions in the literature. In particular: (i) Several papers discuss whether time is an observable in quantum theory. If time were an observable, it should be represented by a self-adjoint operator $`T`$. The spectrum of $`T`$ should be the real line. A well known theorem  demands then its conjugate variable, which is the energy, to have unbounded spectrum. But energy is bounded from below. Therefore time cannot be an observable . But if time is not an observable, how can we measure it? (ii) There are several discussions on observability in general relativity (see  and references therein). In the literature one finds contradictory statements. For instance, that the metric tensor $`g_{\mu \nu }(x)`$ is not observable but a curvature scalar $`R(x)`$ is observable; or that no local quantity such as $`R(x)`$ can be observable. (iii) Observability is a source of lively debates in quantum gravity . Observables must be gauge invariant, therefore commute with the constraints, therefore, in particular, with the Wheeler-DeWitt operator, and therefore they have to be constant in the coordinate time $`t`$. Thus, no quantity that changes with $`t`$ can be observable. This conclusion is considered unreasonable by some . Others (including myself)  argue that the observables in quantum gravity are relative quantities expressing correlations between dynamical variables. But how can a correlation between two non-observable quantities be observable?

I believe that in many debates of this kind there is a confusion between two distinct notions of observability. Mixing up these two notions generates misunderstanding and conceptual mistakes. In this note, I try to clear up the source of this confusion.

The difference between the two notions of observability has to do with localization in time and in space. In a non-relativistic context, the spacetime structure of the world is assumed to be fixed and simple. Because of this, the distinction between the two notions of observability can be disregarded. More precisely, the distinction is replaced by the introduction of a fixed structure on the space of the observables, and then it is safely ignored. The fixed structure of the space of the observables reproduces the fixed structure of spacetime, as we shall see. In a generally relativistic context, on the other hand, the spacetime structure of the world is more complex, and we cannot trade the distinction between different notions of observability for a pre-established structure on the space of the observables. In such a context, ignoring the distinction between different meaning of “observable” leads to serious confusion.

Partial and complete observables are defined in Section II. The two notions are shown to be distinct and examples of the two are given. I then discuss the relevance of the distinction for different contexts: general relativity (in Section III), quantum mechanics (in Section IV), and quantum gravity (in Section V).

The space of the partial observables is the extended configuration space. This space, and its associated extended phase space, on which the hamiltonian constraint is defined, are often presented as devoid of direct physical significance. Instead, I argue in Section VI that the extended configuration space has indeed a direct physical interpretation: it is the space of the partial observables. This space plays a central role in the general structure of mechanics, both at the classical and at the quantum level. I illustrate this role and argue that it provides a unifying perspective that sheds light on the structure of mechanics, especially of general covariant mechanics.

The distinction between partial and complete observables was discussed in . The distinction is sometimes implicitly used, but I am not aware of any other explicit discussion on it in the literature.

## II Partial observables and complete observables

Let us start from the following two definitions***The operational tone of the first definition does not imply any adherence to operationalism here : the reference to measuring procedures is just instrumental for clarifying a distinction.:

a physical quantity to which we can associate a (measuring) procedure leading to a number.

a quantity whose value can be predicted by the theory (in classical theory); or whose probability distribution can be predicted by the theory (in quantum theory).

At first sight, the two definitions might seem equivalent, but they are not. To see this, consider the following example. Imagine we have a bunch of cards in a box. Each card has an upper and a lower side (say, of different colors). On each side, there is a number. Denote the upper number as $`n`$ and the lower number as $`N`$. We extract a certain number of cards from the box and we realize that there is law connecting the two numbers: say $`N`$ is always a certain function of $`n`$. That is $`N=N(n)`$. The law $`N=N(n)`$ gives us a predictive theory for some observable quantities. What are the observables in this context? Clearly, both $`n`$ and $`N`$ are partial observables, according to the definition given above. However, neither of them is predictable, because at each new card we extract we do not know which particular value of $`n`$, or which particular value of $`N`$, will be found. What is predictable is the value of $`N`$ on the back of a card marked with a certain $`n`$. Therefore we have is one “complete observable” $`N(n)`$ for each value of $`n`$. The “complete observables” are $`N(1),N(2),N(3),\mathrm{}`$

The example may seem artificial and unrelated to the structure of realistic physical theories, but it is not. Indeed, realistic physical theories have a structure similar to the one of the example: the role of the “independent” partial observable $`n`$ is played by the quantities giving the temporal localization or the spatio-temporal localization. Consider for instance a very simple physical system, a pendulum. Assume oscillations are small and described by the equation

$$\frac{d^2q(t)}{dt^2}=\omega ^2q(t).$$
(1)
Now suppose we are in a (very simple) laboratory, and we want to check the correctness of (1). What do we need? Clearly we need *two* measuring instruments: one that gives us the pendulum position $`q`$ and one that gives us the time $`t`$. The theory cannot predict the value of $`t`$. Nor can it predict the value of $`q`$, unless we specify that the value of $`q`$ we are interested in is the one at a certain given time $`t`$. Therefore, there are two partial observables playing a role here: $`q`$ and $`t`$. And there is one family of complete observables: the observables $`q(t)`$, for any real value of $`t`$. It is sufficient to know the actual value of a few of these complete observables (for instance $`q(0)`$ and $`dq(t)/dt|_{t=0}`$), in order to be able to predict the value of all the others.

This may seem a rather pedantic account of observability in the context of a non relativistic system. Indeed, one usually says that “$`q`$ is observable”, leaving implicit “yes, of course, one has to say at which time the observation is made”. But, as mentioned, such carelessness in defining observability is then paid for at a high price in a generally relativistic context, where things are not simply evolving in a fixed external time $`t`$ which can be measured by an external clock, as in non-relativistic physics. Let us therefore here clearly distinguish between (i) $`t`$ and $`q`$ (without specified time), which are partial observables, because there are measuring procedures specified for them, but they cannot be predicted, and (ii) the family of complete observables $`q(t)`$, which can be predicted.

Observe that the predictions of a mechanical theory can always be expressed as relations between partial observables. These relations depend on a certain number of parameters, which label the different possible histories of the system. For instance, in the case above the predictions of the theory are given by the following relation between $`t`$ and $`q`$:

$$f(q,t;A,\varphi )=qA\mathrm{sin}(\omega t+\varphi )=0.$$
(2)
From this perspective, the partial observables $`q`$ and $`t`$ can be taken as being on the same footing. That is, they can be treated symmetrically in the theory. Observe, however, that in the example considered the two partial observables $`q`$ and $`t`$ are not entirely on the same footing. The predictions of the theory can certainly be expressed as a relation between the two, but this relation can be solved for $`q`$ as a function of $`t`$, not for $`t`$ as a function of $`q`$. Accordingly, we call $`t`$ an independent partial observable and $`q`$ a dependent partial observable. As we shall see, in a generally relativistic context such distinction between dependent and independent partial observables is lost.

In a non relativistic system with $`m`$ degrees of freedom $`q^i`$, with $`i=1,\mathrm{},m`$, there are in general $`n=m+1`$ partial observables: $`q^a=(t,q^i)`$ with $`a=1,\mathrm{},n`$. The space of these form the extended configuration space of the system, which we denote $`𝒞`$. The predictions of classical mechanics can always be given as relations between the extended configuration space variables, as in (2). These relations depend on a certain number of parameters $`\alpha ^j`$ ($`A`$ and $`\varphi `$ in (2)), which label the different possible histories of the system.

$$f(q^a;\alpha ^j)=0.$$
(3)
Classical mechanics and quantum mechanics can be formulated in a very general and very clean form over the extended configuration space $`𝒞`$. Examples of such formulations are the Hamilton-Jacobi formalism, the extended phase space formalism, the path integral formalism and the propagator formalism . These formulations stress the centrality of the notion of partial observable and show that mechanics treats all partial observables on the same ground. In Section VI, we shall discuss some of these formulations and their relation to partial observability.

Finally, consider a field theory, such as Maxwell electrodynamics. A dynamical variable is represented for instance by the Electric field $`E(\stackrel{}{x},t)`$. The electric field at a given spacetime point $`(\stackrel{}{x},t)`$ can be predicted, and therefore it is a complete observable of the theory. In order to measure $`E(\stackrel{}{x},t)`$, we need *five* partial observables. Indeed, we may imagine that we have at our disposal five measuring devices: a clock measuring $`t`$, an electric field detector that measures $`E`$, and three distance measuring devices, giving the three components of $`\stackrel{}{x}`$. The complete observable $`E(\stackrel{}{x},t)`$ is composed by these five partial observables.

## III General Relativity

Let us now move to a generally relativistic context. For concreteness, let us consider general relativity coupled with $`N`$ small bodies. For instance, these bodies may represent the planets and the satellites in a generally relativistic model of the solar system. The lagrangian variables can be taken to be the metric $`g_{\mu \nu }(\stackrel{}{x},t)`$ and, say, the bodies’ trajectories $`X_{(n)}^\mu (\tau _n)`$, with $`n=1,\mathrm{},N`$ and orientations $`E_{(n)}^{a\mu }(\tau _n)`$ (a local tetrad on the $`n`$th body, $`a=0,\mathrm{},3`$). As well known, the meaning of the coordinates $`(\stackrel{}{x},t)`$ in general relativity is very different from their peaning in pre-general-relativistic (pre-GR) physics.Einstein has described his 1912 to 1915 final struggle for general relativity as a dramatic effort to understand the new “meaning of the coordinates”. Indeed, the coordinates $`(\stackrel{}{x},t)`$ do not represent observable quantities at all.Unless one fixes a physically interpreted gauge, in which case the discussion of observability is a bit different, but the final conclusions are unchanged. See below. That is, the general relativistic coordinates $`\stackrel{}{x}`$ and $`t`$ are neither partial observables nor complete observables.

The distinction between partial and complete observables, however, is still present. Consider some typical predictions of the theory. For instance, a prediction of the theory may be the following: tomorrow morning, when the Sun is 5 degrees over the horizon, Venus will be visible at 12 degrees over the horizon. This is a well defined prediction, and should thus refer to a complete observable. The complete observable is the angle $`\alpha _V`$ that Venus makes with the horizon, at the moment in which the angle $`\alpha _S`$ of the Sun with the horizon is 5 degrees. Clearly, to verify this prediction we need measuring procedures giving us the two angles. Therefore the two angles are partial observables. The complete observable is the value of $`\alpha _V(\alpha _S)`$ for $`\alpha _S=5^o`$.

As a second example, we could replace $`\alpha _S`$ with the proper time $`\tau `$ measured on Earth by a clock that started ($`\tau =0`$) at a certain specified event $`O`$ (say, a certain eclipse). Then again, the proper time $`\tau `$ elapsed from the eclipse, or, equivalently, the length $`\tau `$ of the Earth’s worldline since the eclipses, is a partial observable because it can be measured; but it is not a complete observable, because it cannot be predicted. Indeed, it is an observable quantity used for localizing a spacetime point.

The key difference between general relativity and pre-GR physics as far as observability is concerned is well illustrated by a third example. Consider the following (realistic!) experiment. A very accurate clock is mounted on a satellite. Say a satellite in the GPS system. The satellite broadcasts its local time and the signal is received by the launching base, and compared with the time of an equally accurate clock kept at the base. As well known, the discrepancy between the two due to generally relativistic effects is easily observable using current technology. Let $`\tau _s`$ and $`\tau _b`$ be the signal received from the satellite and the local clock reading. General relativity can be used to predict the relation between the two

$$f(\tau _s,\tau _b)=0,$$
(4)
(once all the relevant initial data are known). Again, we are in a situation of two partial observables forming a complete observable. Now: which one of the two is the independent one? In general, (4) may not be solvable for either variable. One could say that $`\tau _b`$ has to be viewed as the “natural” independent variable, since this is “our” time. But one can equally well say that the $`\tau _s`$ is the “natural” independent variable, since it provides an accepted standard of time … Clearly we are in a very different situation from one with the two partial observables $`q`$ and $`t`$ of the previous section. There, we had a clear distinction between an independent observable ($`t`$) and a dependent one ($`q`$). Here, $`\tau _s`$ and $`\tau _b`$ are truly on the same footing.

The key difference between general relativistic physics and pre-GR physics is the fact that in general relativistic physics the distinction between dependent and independent partial observables is lost. A pre-GR theory is formulated in terms of variables (such as $`q`$) evolving as functions of certain distinguished variables (such as $`t`$). General relativistic systems are formulated in terms of variables (such as $`\tau _b`$, $`\tau _s`$, $`\alpha _V`$, $`\alpha _S`$) that evolve with respect to each other. General relativity expresses relations between these, but in general we cannot solve for one as function of the other. Partial observables are genuinely on the same footing.

What are the complete observables, in general, in this context? A complete observable is a quantity that can be predicted uniquely. Therefore it is a quantity which is well defined once we know the solution of the equations of motion, up to all gauges (that is, which not affected by the indetermination of the evolution). Such a quantity can be seen as a function on the space of the solutions modulo all gauges. This space is the physical phase space of the theory $`\mathrm{\Gamma }`$. In the canonical formalism, $`\mathrm{\Gamma }`$ can be obtained as the space of the orbits generated by the constraints on the constraint surface. Any complete observable can thus be expressed as a function on $`\mathrm{\Gamma }`$. Equivalently, it can be expressed as a function on the extended phase having vanishing Poisson brackets with all first class constraints, including, of course, the hamiltonian constraint. Vice versa, any function that commutes with all constraints defines, in principle, a complete observable.

Partial observables are hard to construct formally in general, but it is far easier to define and use them concretely. For a recent concrete construction of a complete set of partial and complete observables in GR, see .

I close the section with a note on gauge-fixed formulations of GR. One may fix the gauge by choosing coordinates that have a physical interpretation. More precisely, one may select a family of partial observables (curvature scalars, scalar fields, dust variables, GPS readings…) and fix the coordinate gauge by tying the coordinate system to these partial observables. Within a formulation of this kind, coordinates represent partial observables. Furthermore, they have a natural character of independent partial observables. However, this does imply that the independent partial observables are determined by the theory, because the same physical situation can be described by a different physical gauge choice, in which the role of dependent and independent partial observables is interchanged.

## IV Quantum theory

In quantum theory observables are represented by operators. Which observables are represented by operators: the partial or the complete observables? The answer is different in the Heisenberg picture (evolving operators) and in the Schrödinger picture (evolving states). Let us start from the Heisenberg picture. Here the operators are time dependent. For instance, in the quantum theory of a harmonic oscillator in the Heisenberg picture, there is no position operator $`Q`$, but only the operator $`Q(t)`$ that represents “position at time $`t`$”. This is immediately recognized as the operator corresponding to the complete observable $`q(t)`$ discussed in Section II. In the Heisenberg picture operators are associated with complete observables.

In the Schrödinger picture, there is an operator $`Q`$ associated with the partial observable $`q`$. However, specific predictions are not given just in terms of this operator: we need the state as well, and, in the Schrödinger picture, the state $`\mathrm{\Psi }(t)`$ is time dependent. Thus, for instance, the expectation value $`\overline{Q}(t)=\mathrm{\Psi }(t)Q\mathrm{\Psi }(t)`$, which is a prediction of the theory, is associated with the complete observable $`q(t)`$, as it should, not with the partial observable $`q`$. In order for the Schrödinger picture to be meaningful, we need the theory to be expressed in terms of a well defined independent partial observable $`t`$: “the external time”. In a theory such as general relativity, where the dynamics expresses the relation between partial observables that are on equal footing, the Schrödinger picture is not viable. More precisely, it will be viable only in special circumstances, in which we can choose (arbitrarily) one of the partial observables as the independent one and solve the dynamical relations expressing the predictions of the theory in terms of this quantity. In general, no such quantity exists. On the other hand, the Heisenberg picture remains meaningful whatever the spacetime structure of the theory. Let us therefore return to the Heisenberg picture, which is far more general.

In the Heisenberg context, consider the problem of whether there should be a time operator in quantum theory. The time $`t`$ is a partial observable, not a complete observable. Operators are associated with complete observables, not with partial observables. Therefore it is against the tenets of quantum theory to search for an operator corresponding to $`t`$. Operators correspond to quantities that are in principle predictable (such as $`q(t)`$), not to quantities (such as $`t`$) that serve only to localize the measurement of a predictable quantity in spacetime. A quantity that is described by an operator in quantum theory is a quantity such that there are states that diagonalize it, namely such that there are physical situations in which the outcome of a measurement of that quantity is certain: the time $`t`$, on the contrary, can never be predicted.

In other words, quantum theory deals with the relation between $`q`$ and $`t`$, and not with $`q`$ alone or $`t`$ alone. Therefore it is meaningless to search for the quantum theory of the $`t`$ variable alone.

Of course, the reading $`T`$ of a clock can be predicted, but only if we first read another clock. If we know that the second clock indicates $`t`$, we can predict that the first clock will read a certain $`T`$. If we now take into account the fact that the clock is a physical mechanical system and it is subject to quantum fluctuations, then we can describe it in terms of an operator. This operator describes the complete observable $`T(t)`$. There will be quantum fluctuations described by generic states in the state space on which this operator acts. These fluctuations are not the quantum fluctuations of one independent time variable. They are the quantum fluctuations in the observable correlation between two clock variables.

## V Quantum gravity

In quantum gravity, operators corresponding to physical observables must commute with the Wheeler-DeWitt constraint operator. This operator is the generator of evolution in the coordinate time $`t`$. Thus, physical observables must be invariant under evolution in $`t`$. This fact has raised much confusion. How can observables invariant under evolution in $`t`$ describe the evolution we observe? The question is ill posed, because it confuses evolution with respect to the coordinate time $`t`$ and physical evolution. In Section III we have observed that in general relativity quantities like the proper times $`\tau _b`$ and $`\tau _s`$ are partial observables and their relative evolution is well defined. Let us fix a value $`\tau _b=\tau `$ of the first, and consider the corresponding value of $`\tau _s`$. (If there are several such values, take the highest). Call this value $`T_\tau `$. That is $`T_\tau `$ is the highest number for which

$$f(T_\tau ,\tau )=0.$$
(5)
where $`f`$ is the function in (4). $`T_\tau `$ is a complete observable. It is the signal we receive from the satellite when our local proper time at the base is $`\tau `$. It describes the change of the value of the received signal as the proper time at the base passes. This is a description of evolution. At the same time, this is a quantity independent of the coordinate $`t`$. To see this, recall that to calculate its value from a specific solution of the Einstein equations, we first find the dependence of $`\tau _b`$ and $`\tau _s`$ on the coordinate time $`t`$. Namely we compute the functions $`\tau _b(t)`$ and $`\tau _s(t)`$. The form of these two functions is gauge dependent: it changes if we use a different coordinate representation of the same four geometry. We then locally invert the second function and insert $`t(\tau _s)`$ in the first. The resulting $`T_\tau \tau _s(\tau _b=\tau )`$ is independent of the coordinate $`t`$ chosen, and thus it is uniquely determined by the equivalence class of solutions of the field equations under diffeomorphisms. It is a well defined on the space of these equivalence classes, namely on $`\mathrm{\Gamma }`$. Equivalently, it can be represented as a function on the extended phase space that commutes with all the constraints, including the hamiltonian constraint.

Let us now come to a main objection that we want to address in this paper, which is the following.

$`T_\tau `$ cannot be observable without $`\tau _b`$ and $`\tau _s`$ being individually observable. Thus $`\tau _b`$ and $`\tau _s`$ are observable. Observables must be represented by physical operators. $`\tau _b`$ and $`\tau _s`$ depend on $`t`$ and do not commute with the hamiltonian constraint. Therefore in any quantum theory of gravity there should be physical operators representing observables that do not commute with the Wheeler-DeWitt constraint.

It should be clear at this point why this objection is wrong. It confuses partial and complete observables. $`\tau _b`$ and $`\tau _s`$ are partial observables, and partial observables are not associated with quantum operators in quantum theory (more precisely, in Heisenberg picture quantum theory, which is the only one viable in this context).

We close this section with an observation on the role of the coordinates in the formalism of quantum gravity. The general relativistic spacetime coordinates $`(\stackrel{}{x},t)`$ have no direct physical interpretation. In a gauge fixed context, they can be tied to partial observables. In any case, however, they do not represent complete observables. It follows that the idea that the coordinates should be represented by quantum operators is not justified is the light of quantum theory and general relativity alone. Operators are attached to complete observables, while spacetime coordinates are –at best– partial observables: they cannot be predicted, they serve only to localize complete observables.

Quantum theory deals with the relation between partial observables. It can deal with the relation between physical variables and (gauge-fixed) coordinates $`(\stackrel{}{x},t)`$. But not with the value of the coordinates alone. Therefore it is meaningless to search for the quantum theory of the $`(\stackrel{}{x},t)`$ variables alone.

Non-commutative geometry approaches to quantum gravity search for a mathematics capable of promoting the spacetime coordinates $`(\stackrel{}{x},t)`$ to a non-commuting operator algebra. This approach is sometimes motivated with the argument that quantum theory should require the coordinates $`(\stackrel{}{x},t)`$ to be represented by operators. In the light of the discussion above, I think that this motivation mistakes partial observables and complete observables. Non-commutative approaches to quantum gravity are extremely interesting in my view, both mathematically and the physically. But I think that this particular motivation is naive and not tenable. Physical non-commutativity of quantities related to physical localization and geometry, on the other hand, is likely to follow from the fact such quantities should, in fact, be functions of the gravitational field, and therefore quantum dynamical variables.

## VI The extended configuration space

In this section I discuss the role of the partial observables in the formal structure of mechanics. I focus here on theories with a finite number of degrees of freedom, leaving the extension to field theory to the reader.§§§Examples of generally relativistic systems with a finite number of degrees of freedom are provided for instance by cosmological models, or by models as in . This discussion sheds light on the physical interpretation of certain structures, such as the extended phase space of the fully constrained systems. For a more complete treatement of this subject, see ; see also . The central message of this section is double. First, that when the formalism is sufficiently general, partial observables are the main quantities mechanics deals with. Second, that, in general, mechanics makes no distinction between dependent and independent observables. The distinction between independent and dependent observables can be seen as an accident of the specific dynamics of non-relativistic theories. In the light of these two observations, I think that the interpretation of general relativistic theories becomes more transparent.

As observed in Section II, the partial observables of a mechanical system form the extended configuration space $`𝒞`$. Recall that we denote the partial observables by $`q^a`$, $`a=1,\mathrm{},n`$. Dynamics can be given in terms of a first order partial differential equation on $`𝒞`$, the Hamilton-Jacobi equation

$$C(q^a,\frac{S(q^a)}{q^a})=0,$$
(6)
The function of 2$`n`$ variables $`C(q^a,p_a)`$ determines the dynamics. One searches for an $`n`$-parameter family of solutions of this equation $`S(q^a,Q^a)`$, where $`Q^a`$ are $`n`$ constants, and the predictions of the theory are contained in equations (3), which are obtained as follows

$$f_a(q^a;Q^a,P_a)=\frac{S(q^a,Q^a)}{Q^a}P_a=0.$$
(7)
These form a $`2n`$-parameter family of $`n`$ relations between the partial observables (not all independent). The parameters $`Q^a,P_a`$ label the possible histories of the system: each history determines a set of relations among partial observables. These relations are the physical predictions of the theory. Notice that all partial observables $`q^a`$ are treated on the same footing: the “time” partial observable, if present at all, is just a variable among the others. Notice also that the usual last step of the Hamilton-Jacobi prescription, which is to invert (7) for the dependent variables, is not necessary from this point of view.

Let $`\mathrm{\Gamma }`$ be the space of the histories. Generically, there is one history connecting any two points of $`𝒞`$. Therefore $`\mathrm{\Gamma }`$ has dimension $`2n2`$. Since $`Q^a`$ and $`P_a`$ are $`2n`$ functions on $`\mathrm{\Gamma }`$, they over-coordinatize $`\mathrm{\Gamma }`$ and there are 2 relations among them. Also, histories are one-dimensional, and therefore only $`n`$-1 of the $`n`$ relations (7) are independent. The space $`\mathrm{\Gamma }`$ is the phase space of the system. A point in $`\mathrm{\Gamma }`$ is a “state” of the theory, in the sense of “Heisenberg state‘” . It represents a possible history of the system, not a “state at a certain time”.Dirac argued repeatedly that the Heisenberg notion of state is the good one, and the only one that makes sense in a relativistic context. See for instance Sec I.3 of the first edition of . In later editions of this book Dirac shifted the emphasis to the Schrödinger states, explaining (in the Preface) that these, after all, are easier to work with in the non relativistic context, although “it seems a pity” to give the cleaner notion.

The function $`f`$ in (7) is defined on the cartesian product of the space of the partial observables with the space of the states.

$$f:𝒞\times \mathrm{\Gamma }IR^n.$$
(8)
The entire predictive content of a dynamical theory is in the surface $`f=0`$ in the cartesian product of the space of the partial observables and the space of the states. For each point $`q^a`$ in $`𝒞`$, the surface $`f=0`$ determines the set of states compatible with the value $`q^a`$ of the partial observables. For each state in $`\mathrm{\Gamma }`$, the surface determines a relation among the partial observables in $`𝒞`$.

In the special case of a non-relativistic system, one of the partial observables $`q^a`$ is the time $`t`$. Let it be, say, $`q^0`$, and call the other partial observables $`q^i`$ with $`i=1,\mathrm{},m=n1`$. In this case the function $`C(q^a,p_a)`$ has the special form

$$C(q^a,p_a)=p_0+H(q^i,p_i).$$
(9)
Therefore in this special case the Hamilton-Jacobi equation takes the well known form

$$\frac{S(q^i,t)}{t}+H(q^i,\frac{S(q^i,t)}{q^i})=0.$$
(10)
The general Hamilton-Jacobi formalism has a nice geometrical interpretation in the canonical framework. Let us illustrate it, with the purpose of discussing the meaning of the structures of generally covariant hamiltonian systems.

Consider the cotangent bundle $`T^{}𝒞`$ over the extended configuration space, with canonical coordinates $`(q^a,p_a)`$. Call it the extended phase space. It carries the natural Poincaré one-form $`\theta =p_adq^a`$, and the symplectic form $`\omega =d\theta `$. The dynamics is coded in a relation on $`T^{}𝒞`$:

$$C(q^a,p_a)=0,$$
(11)
In the special case of a non-relativistic system, $`q^a=(q^0,q^i)`$ and Equation (11) has the form

$$C(q^a,p_a)=p_0+H(q^i,p_i)=0$$
(12)
where $`H`$ is the Hamiltonian. The variable $`q^0=t`$ is the time variable, and its conjugate momentum $`p_0=E`$ is (minus) the energy. The dynamics of the system is then coded in the relation (12) which gives the energy as a function of the other coordinates and momenta.

Equation (11) defines a surface $`\mathrm{\Sigma }`$ in $`T^{}𝒞`$. Call $`\omega _\mathrm{\Sigma }`$ the restriction of $`\omega `$ to this surface. The “presymplectic” two-form $`\omega _\mathrm{\Sigma }`$ is degenerate and has a null direction. It is not difficult to see that the integral curves of this null direction are the solutions of the equations of motion of the system.The coordinate form of the relation $`Y(\omega _\mathrm{\Sigma })=0`$ between $`\omega _\mathrm{\Sigma }`$ and its null vector field $`Y`$ is given by the Hamilton equations. The space of these curves is the physical phase space of the system $`\mathrm{\Gamma }`$ and carries a unique symplectic two-form $`\omega _\mathrm{\Gamma }`$ whose pull back to $`\mathrm{\Sigma }`$ under the natural projection $`\pi :\mathrm{\Sigma }\mathrm{\Gamma }`$ is $`\omega _\mathrm{\Sigma }`$. Let $`P_a`$ and $`Q^a`$ be coordinates that (over-)coordinatize\**\**\**$`2(n1)`$ coordinates are sufficient to coordinatize $`\mathrm{\Gamma }`$. For instance, one can take initial coordinate and momenta at $`t=t_0`$. We prefer to use here $`2n`$ coordinates for reasons that will be clear below. The extra coordinates can be seen as the initial time $`t=t_0`$ and the energy. A change in the first amounts to a relabeling of the meaning of the initial data. The second is constrained by $`C`$. $`\mathrm{\Gamma }`$ and define a one-form $`\theta _\mathrm{\Gamma }=P_adQ^a`$ such that $`d\theta _\mathrm{\Gamma }=\omega _\mathrm{\Gamma }`$. Then (using the same notations for the forms and their pull back), since on $`\mathrm{\Sigma }`$ we have $`\omega _\mathrm{\Gamma }=\omega _\mathrm{\Sigma }`$, it follows that

$$d(p_adq^a)=d(P_adQ^a)$$
(13)
or

$$p_adq^a=P_adQ^a+dS$$
(14)
where $`S`$ is a zero-form on $`\mathrm{\Sigma }`$. But let us pull the coordinates $`Q^a`$ back onto $`\mathrm{\Sigma }`$ and assume that the set $`(q^a,Q^a)`$ (over-) coordinatizes $`\mathrm{\Sigma }`$. Then (14) gives

$`{\displaystyle \frac{S(q^a,Q^a)}{q^a}}`$ $`=`$ $`p_a,`$ (15)
$`{\displaystyle \frac{S(q^a,Q^a)}{Q^a}}`$ $`=`$ $`P_a.`$ (16)
From (15) and (11) we obtain

$$C(q^a,\frac{S(q^a,Q^i)}{q^a})=0,$$
(17)
which is the Hamilton-Jacobi equation (6) and can be used to compute $`S`$; while (16) is the equation (7) giving the physical predictions from $`S`$.

What is the physical meaning of $`S(q^a,Q^a)`$? Without loss of generality, we can choose the integration constants so that

$$S(q^a,Q^a=q^a)=0.$$
(18)
Fix a point $`p`$ on $`\mathrm{\Sigma }`$, and consider the trajectory that starts on $`p`$. Along this trajectory $`dQ^a=0`$ and thus from (14) we have $`dS=p_adq^a`$. Parametrize the trajectory with an arbitrary time parameter $`t`$ and write $`dS=p_adq^a=p_a\dot{q}^adt`$. The canonical hamiltonian with respect to this parameter is null, and therefore $`p_a\dot{q}^aL=0`$. Therefore we have $`dS=Ldt`$ along each orbit $`q^a(t)`$ and

$$S(q^a,Q^a)=_{Q^a}^{q^a}L(q^a(t))𝑑t.$$
(19)
That is, $`S(q^a,Q^a)`$ is the action, computed over the physical trajectory that joins the points with coordinates $`Q^a`$ and $`q^a`$. In the case of a non-relativistic system, let $`q^a=(q^i,t)`$. Then $`dS=p_adq^a=p_idq^iHdt`$. Recall that $`H=p_i\dot{q}^iL`$, where $`L`$ is the Lagrangian. Therefore $`dS=p_idq^ip_idq^i+Ldt=Ldt`$ along each orbit. Thus,

$$S(q^a,Q^a)=_{Q^i}^{q^i}L(q^i(t))𝑑t.$$
(20)
where the trajectory starts at time $`t`$ in $`q^i`$ and ends at time $`T`$ in $`Q^i`$. That is, $`S(q^a,Q^a)`$ is still the action, computed over the physical trajectory that joins the points with coordinates $`Q^a`$ and $`q^a`$.

Notice that from this point of view Hamilton’s principal function and Hamilton’s characteristic function are identified. More precisely, $`S(q^a,Q^a)`$ is the principal function for the evolution in any partial observable identified as the time $`q^0=t`$. But it is also the characteristic function of the evolution in an arbitrary parameter time along the histories. And it is also the principal function for the evolution in such a time, since the hamiltonian that generates this motion vanishes. This compactification of the formalism is quite remarkable.

In conclusion, the ingredients of mechanics can be taken to be solely the extended configuration space $`𝒞`$ and the function $`C`$ on $`T^{}𝒞`$. A mechanical system is determined by the pair $`(𝒞,C)`$. The kinematics of a specific theory is determined by the space of its partial observables $`𝒞`$; its dynamics is determined by the constraint $`C(q^a,p_a)=0`$ on the associated phase space. There is no need to single out a specific partial variable as the time, nor to mention evolution. Mechanics is a theory of relations between partial observables. No distinction between dependent and independent partial observables is required. This distinction is an accident of non-relativistic theories, in which the constraint $`C(q^a,p_a)`$ happens to have the form (12).

Why do I stress this fact? Because generally relativistic theories are formulated in terms of constraints such as (11) over an extended configuration space. It is sometimes claimed that the theory can only be interpreted if one finds a way to “deparametrize” the theory, namely to select the independent variable, among the variables $`q^a`$. In the opposite camp, the statement is sometimes made that only variables on the physical phase space $`\mathrm{\Gamma }`$ have physical interpretation, and no interpretation should be associated with the variables of the extended configuration space $`𝒞`$. Instead, I have argued here that the variables of the extended configuration space have a physical interpretation as partial observables. In a sense, they are the quantities with the most direct physical interpretation in the theory.

Finally, consider quantum theory. The Schrödinger equation, as well as the Wheeler-DeWitt equation, are both partial differential equations on the extended configuration space $`𝒞`$. They can both be obtained in general from the constraint (11), with no need of distinguishing dependent from independent partial observables. Indeed, they are obtained as

$$C(q^a,i\mathrm{}\frac{}{q^a})\psi (q^a)=0.$$
(21)
The usual physical scalar product on an appropriate space of the solutions of this equation has an intrinsic meaning and does not need the time variable to be singled out in order to be defined – see for instance . All relevant physical predictions of the theory can be extracted from the knowledge of the propagator $`W(q^a,Q^a)`$, which satisfies

$$C(q^a,i\mathrm{}\frac{}{q^a})W(q^a,Q^a)=0.$$
(22)
The propagator gives the probability amplitude for finding the combination of partial observables $`q^a`$ if the combination of partial observables $`Q^a`$ was previously observed. Virtually all predictions of quantum mechanics can be formulated in this covariant manner, on the extended configuration space. This is discussed in detail in . As well known, in the limit of small $`\mathrm{}`$ the Schrödinger equation (or the Wheeler-DeWitt equation) goes over into the Hamilton-Jacobi equation, and the propagator $`W(q^a,Q^a)`$ is given to first order just by the exponential of the action $`W(q^a,Q^a)\mathrm{exp}\{iS(q^a,Q^a)/\mathrm{}\}`$.

## VII Conclusions

I have observed that the notion of observable is ambiguous, and I have discussed the distinction between partial observables and complete observables. This distinction clarifies a certain number of issues related to observability. In particular, I have examined the role played by this distinction in general relativity, in quantum mechanics and in quantum gravity.

The partial observables form the extended configuration space $`𝒞`$. This space seems to be a natural home for classical and quantum mechanics. The two theories admit a clean formulation over this space, which is sufficiently general to deal naturally with general relativistic systems.

A mechanical system is a pair $`(𝒞,C)`$. The space of the partial observables $`𝒞`$ describes the kinematics of the theory. $`C`$ is a function on $`T^{}𝒞`$ that determines the dynamics. Classical dynamics is about relations between partial observables. These relations depend on a certain number of parameters, which label the (time independent) states of the system. The space of these states is the phase space $`\mathrm{\Gamma }`$. The predictions of the theory are therefore given by a surface $`f=0`$ on $`𝒞\times \mathrm{\Gamma }`$. The surface $`f=0`$, as well as $`\mathrm{\Gamma }`$, are determined by the pair $`(𝒞,C)`$. (The general structure of classical and quantum mechanics in this language is discussed more in detail in .)

By fixing a subset of partial observables (one for a mechanical system, four for a field theory), the other partial observables are determined as functions on $`\mathrm{\Gamma }`$. This defines the complete observables of the theory, whose value can be predicted uniquely if the state is known.

Quantum mechanics gives the probability amplitude $`W(q^a,Q^a)`$ for measuring the combination of partial observables $`q^a`$ after having measured the combination $`Q^a`$. Alternatively, it gives the probability distribution for the different possible outcomes of a measurement of the complete observables. These are represented by self-adjoint operator over the (Heisenberg) state space.

No distinction between independent and dependent partial observable is required. The different partial observables can be viewed as being on the same footing. This formulation of mechanics does not require the notion of external time. It is therefore appropriate for general relativistic systems, which are not formulated in terms of evolution in time.

——————

I thank Julian Barbour for extensive discussions and for help with the language. This work was partially supported by NSF Grant PHY-9900791.


text:
# Chemical potential response of pseudoscalar meson masses in the Nambu–Jona-Lasinio model In memory of Prof. Osamu Miyamura, a very kind and warm-hearted person as well as a great physicist.

## 1 Introduction

While the structure of QCD at finite temperature has been investigated in detail using lattice QCD, little is known about matter at finite baryon density due to the sign problem . In order to avoid this difficulty one can consider 2–color QCD , imaginary chemical potential simulations , the response of a hadron mass to changes in the chemical potential, $`\frac{m}{\mu }`$ and $`\frac{^2m}{\mu ^2}`$, at zero chemical potential ($`\mu `$=0) , and so on.

Since the direct application of QCD at finite temperature and chemical potential is not available in the present lattice QCD simulations, effective models of QCD are commonly used. One of the most popular models is the Nambu–Jona-Lasinio (NJL) model . This model has been widely used for describing hadron properties in hot and/or dense matter .

In this work we present NJL model calculations of $`\frac{m}{\mu }`$ for the pion and kaon. The primary goal of our study is to get the same quantity which is simulated on the lattice, i.e., $`\frac{m}{\mu }`$($`T`$, $`\mu `$=0). Of course, the direct comparison between the lattice data and the present NJL model calculations is rather difficult because $`\frac{m}{\mu }`$ on the lattice means the variation of a screening mass, while it corresponds to the variation of a dynamical mass in the present work. Nevertheless, we can compare both results to each other qualitatively.

In contrast to the lattice simulations, we also get $`\frac{m}{\mu }`$ at finite chemical potential within this model. Then, this may give information on the role of the light quark (the $`u`$ and $`d`$ quarks) chemical potential and/or the strange quark chemical potential in hot and/or dense matter. In addition, we consider two cases for the pion. One is the pion with the same $`u`$, $`d`$ quark mass. The other is the pion with non-degenerate $`u`$ and $`d`$ quark masses. Since the present lattice QCD can not simulate the characteristics coming from slightly different $`u`$ and $`d`$ quark masses , the response of the pion mass with different $`u`$ and $`d`$ quark masses can be useful for future studies of $`\frac{m}{\mu }`$ on the lattice.

Due to the limitation of space we present only the results with a non-zero vector–axial vector coupling $`g_V`$ in this paper. Please refer to  for basic formulae to get $`\frac{m}{\mu }`$ for the kaon and pion in the NJL model and the results without $`g_V`$. Detailed derivations including small corrections in the figures in  will be given elsewhere .

## 2 The effects of a non-zero $`g_V`$

One of NJL model calculations showed that the $`K^{}`$ mass at finite density with $`g_V`$0 is quite different from that with $`g_V`$=0 . In this section, we show the chemical potential response of the kaon mass by including a non-zero $`g_V`$. In the case of the pion, the effect of the non-zero vector–axial vector coupling is negligible .

We follow the formalism and the same parameters in . First, in Fig. 1 let us show the responses of quark condensates with (and also without) $`g_V`$ at $`\mu _u=\mu _d`$=0.02, 0.04 GeV, respectively. The effect of a non-zero $`g_V`$ on the $`s`$ quark condensate is much larger than that on the $`u`$ quark condensate. This behavior can be understood from a large difference in the $`s`$ quark mass, i.e., $`m_s`$=135.7 MeV in  and $`m_s`$=88.0 MeV in the present work. Since the response of the $`s`$ quark condensate becomes much larger than before, one can expect that $`\frac{m_K}{\mu }`$ for the kaon will be quite different from the previous values.

In Fig. 2, we show $`\frac{m_K}{\mu _S}`$ for $`K^+`$ and $`K^{}`$. For comparison we also show the previous results without $`g_V`$. The comparison of $`\frac{m_K}{\mu _S}`$ at each point is rather meaningless because all the parameters including the quark masses are changed. We would like to show only a qualitative behavior of $`\frac{m_K}{\mu _S}`$ in this figure. The mass of $`K^+`$ is slightly affected by the vector-axial vector interaction. On the other hand, the behavior of $`\frac{m_K}{\mu _S}`$ for $`K^{}`$ is much different from the previous one, where $`g_V`$=0. This result confirms the previous NJL model calculations that the vector–axial vector interaction reduces the effects of the Fermi sea and the $`K^{}`$ mass is a smoothly decreasing function of density .

## 3 Concluding Remarks

Using the NJL model we have calculated the chemical potential responses of the kaon and pion masses, $`\frac{m_K}{\mu }`$ and $`\frac{m_\pi }{\mu }`$, at zero and finite chemical potential, and found that their behaviors are quite different from each other. Our results show that $`\frac{m}{\mu }`$ is much dependent on the mass difference of two quarks, i.e., the mass difference between the $`u`$ and $`s`$ quarks, or even between the $`u`$ and $`d`$ quarks. Even at very small chemical potentials $`\frac{m_K}{\mu }`$ is much different from $`\frac{m_\pi }{\mu }`$.

Let us discuss some uncertainties in our calculations. First, the mass in our formalism means the dynamical mass, while the screening mass is simulated on the lattice. For the direct comparison we may need the formalism to calculate the screening mass in the NJL model as given in . This may help us to properly interpret the lattice data. Second, while the lattice simulations of $`\frac{m}{\mu }`$ show a large difference between in the confinement phase and in the deconfinement phase , we can not predict $`\frac{m}{\mu }`$ above the Mott temperature in the present work. This is because we have excluded the imaginary part of the 1–loop polarization function in the dispersion relation and considered only the region below the Mott temperature . Thus, the imaginary part of the polarization function should be included for a detailed analysis above the Mott temperature.

As a final remark, we confirm that $`\frac{^2|\overline{q}q|}{\mu ^2}`$ is negative within the present NJL model approach, and find a clear quark mass dependence between $`\frac{^2|\overline{u}u|}{\mu ^2}`$ and $`\frac{^2|\overline{d}d|}{\mu ^2}`$ both below and above the Mott temperature . The second order response of the quark condensate $`\frac{^2\overline{q}q}{\mu ^2}`$ at $`\mu `$=0 can also be simulated on the lattice , and its behavior is similar to the result from the present NJL model calculation.

## Acknowledgements

We thank A. Nakamura for valuable comments. The work of S.C. was supported by the Japan Society for the Promotion of Science (JSPS).


text:
# The Number of Point-Splitting Circles

## 1 Introduction

Our starting point is the following problem, which first appeared in the 1962 Chinese Mathematical Olympiad .

Problem 1.1. Let $`S`$ be a set of $`2n+1`$ points in the plane such that no three are collinear and no four are concyclic. Prove that there exists a circle which has 3 points of $`S`$ on its circumference, $`n1`$ points in its interior, and $`n1`$ in its exterior.

Following \[3, p.48\], we call such a circle point-splitting for the given set of points. For the rest of sections 1 and 2, $`S`$ denotes an arbitrary set of $`2n+1`$ points in general position in the plane, where $`n`$ is a fixed integer.

There are several solutions to problem 1.1. Perhaps the easiest one is the following. Let $`A`$ and $`B`$ be two consecutive vertices of the convex hull of $`S`$. We claim that some circle going through $`A`$ and $`B`$ is point-splitting. All circles through $`A`$ and $`B`$ have their centers on the perpendicular bisector $`\mathrm{}`$ of the segment $`AB`$. Pick a point $`O`$ on $`\mathrm{}`$ which lies on the same side of $`AB`$ as $`S`$, and is so far away from $`AB`$ that the circle $`\mathrm{\Gamma }`$ with center $`O`$ and going through $`A`$ and $`B`$ completely contains $`S`$. This can clearly be done. Now slowly “push” $`O`$ along $`\mathrm{}`$, moving it towards $`AB`$. The circle $`\mathrm{\Gamma }`$ will change continuously with $`O`$. As we do this, $`\mathrm{\Gamma }`$ will stop containing some points of $`S`$. In fact, it will lose the points of $`S`$ one at a time: if it lost $`P`$ and $`Q`$ simultaneously, then points $`P,Q,A`$ and $`B`$ would be concyclic. We can move $`O`$ so far away past $`AB`$ that, in the end, the circle will not contain any points of $`S`$.

Originally, $`\mathrm{\Gamma }`$ contained all the points of $`S`$. Now, as it loses one point of $`S`$ at a time, we can decide how many points we want it to contain. In particular, if we stop moving $`O`$ when the circle is about to lose the $`n`$-th point $`P`$ of $`S`$, then the resulting $`\mathrm{\Gamma }`$ will be point-splitting: it will have $`A`$, $`B`$ and $`P`$ on its circumference, $`n1`$ points inside it, and $`n1`$ outside it, as shown in Figure 1.

The above proof hints that any set $`S`$ has several different point-splitting circles. We can certainly construct one for each pair of consecutive vertices of the convex hull of $`S`$. In fact, the argument above can be modified to show that, for any two points of $`S`$ we can find a point-splitting circle going through them. The reader might find it instructive to work out a proof.

This suggests that we ask the following question. What can we say about the number $`N_S`$ of point-splitting circles of $`S`$? At first sight, it seems that we really cannot say very much about this number. Point-splitting circles seem hard to “control”, and harder to count.

We should be able to find upper and lower bounds for $`N_S`$ in terms of $`n`$. Right away we know that $`N_Sn(2n+1)/3`$, since we can find a point-splitting circle for each pair of points of $`S`$, and each such circle is counted by three different pairs. Computing an upper bound seems more difficult. If we fix points $`A`$ and $`B`$ of $`S`$, it is indeed possible that all $`2n1`$ circles through $`A`$, $`B`$ and another point of $`S`$ are point-splitting. The reader is invited to check this. This is not likely to happen very often in a set $`S`$, and we can get some upper bound out of this. However, it is hard to make this precise and get a non-trivial upper bound.

When $`S`$ consists of 5 points, the situation is simple enough that we can actually show that $`N_S=4`$ always. This was done in . It was also proposed, but not chosen, as a problem for the 1999 International Mathematical Olympiad. Notice that our lower bound above gives $`N_S4`$.

In a different direction, problem 5 of the 1998 Asian-Pacific Mathematical Olympiad, proposed by the author, stated the following.

Proposition 1.2. $`N_S`$ has the same parity as $`n`$.

This result follows easily from the nontrivial observation that, for any $`A`$ and $`B`$ in $`S`$, the number of point-splitting circles that go through $`A`$ and $`B`$ is odd.

The following result brings together the above considerations.

Theorem 1.3. Any set $`S`$ of $`2n+1`$ points in the plane in general position has exactly $`n^2`$ point-splitting circles.

Theorem 1.3 is the main result of this paper. In section 2 we prove that every set of $`2n+1`$ points in the plane in general position has the same number of point-splitting circles. In section 3 we prove that this number is exactly $`n^2`$. In section 4 we present some questions that arise from our work.

## 2 $`N_S`$ is Constant

At this point, we could go ahead and prove the very counterintuitive Theorem 1.3, suppressing the motivation behind its discovery. With the risk of making the argument seem longer, we believe that it is worthwhile to present a natural way of realizing and proving that the number of point-splitting circles of $`S`$ depends only on $`n`$. Therefore, we ask the reader to forget momentarily the punchline of this article.

Suppose that we are trying to find out whatever we can about the number $`N_S`$. As mentioned in Section 1, this number does not seem very tractable and it is not clear how much we can say about it. Being optimistic, we can hope to be able to answer the following two questions.

Question 2.1. What are the sharp lower and upper bounds $`m=m_{2n+1}`$ and $`M=M_{2n+1}`$ for $`N_S`$?

Question 2.2. What are all the values that $`N_S`$ takes in the interval $`[m,M]`$?

Question 2.1 seems considerably difficult. To answer it completely, we would first need to prove an inequality $`mN_SM`$, and then construct suitable sets $`S_{min}`$ and $`S_{max}`$ which achieve these bounds. To see how difficult this is, the reader is invited to try to construct any set $`S`$ of $`2n+1`$ points for which the number $`N_S`$ can be easily computed.

At this point question 2.1 seems very hard, so let us focus on Question 2.2 instead. Here is a first approach.

Intuitively, since the set $`S`$ can be transformed continuously, we should expect the value of $`N_S`$ to change “continuously” with it. Suppose we start with the set $`S_{min}`$ (with $`N_S=m`$) and move its points continuously so that we end up with $`S_{max}`$ (with $`N_S=M`$). The value of $`N_S`$ should change “continuously” as $`S`$ changes continuously. By “continuity” we would guess that $`N_S`$ sweeps all the integers between $`m`$ and $`M`$ as $`S`$ changes from $`S_{min}`$ to $`S_{max}`$.

Right away, we know that this is not entirely true. By Proposition 1.2 we know that the parity of $`N_S`$ is determined by $`n`$, so $`N_S`$ will not sweep all the integers between $`m`$ and $`M`$. This is not too surprising, since we haven’t made precise the meaning of the statement that the value of $`N_S`$ should change “continuously” as $`S`$ changes continuously. The above guess assumed that the value of $`N_S`$ can only jump by 1 as $`S`$ is transformed continuously. (That is, if we have a set with $`k1`$ point-splitting circles and we deform it continuously into a set with $`k+1`$ point-splitting circles, then somewhere in the middle we must have had a set with $`k`$ point-splitting circles.) We have no reason to assume that.

We can still hope that, as $`S`$ changes, $`N_S`$ sweeps all the integers of the right parity between $`m`$ and $`M`$. To show this, we would have to show that the value of $`N_S`$ can only jump by 2 as $`S`$ is transformed continuously. This is a reasonable statement which we can try to prove.

In any case, the natural question to ask is what kind of “continuity” the value of $`N_S`$ satisfies as $`S`$ changes continuously. We certainly expect that if two sets $`S`$ and $`T`$ look very very much alike, then the difference $`N_SN_T`$ should be small. We have to find a way to make this statement precise.

Suppose we have sets $`S_{min}=\{P_1,\mathrm{},P_{2n+1}\}`$ and $`S_{max}=\{Q_1,\mathrm{},Q_{2n+1}\}`$ that achieve the upper and lower bounds for $`N_S`$, respectively. Now slowly transform $`S_{min}`$ into $`S_{max}`$: first send $`P_1`$ to $`Q_1`$ continously along some path, then send $`P_2`$ to $`Q_2`$ continuously along some other path, and so on. We can think of our set $`S`$ as changing with time. At the initial time $`t=0`$, our set is $`S(0)=S_{min}`$. At the final time $`t=T`$, our set is $`S(T)=S_{max}`$. In between, $`S(t)`$ varies continously with respect to $`t`$. How does $`N_{S(t)}`$ vary “continuously” with time? How small can we make $`N_{S(t+\mathrm{\Delta }t)}N_{S(t)}`$ for small enough $`\mathrm{\Delta }t`$? This is the question we need to ask.

Technical Remark. As we move from $`S(0)`$ to $`S(T)`$ continuously, it is likely that several intermediate sets $`S(t)`$, with $`0<t<T`$, are not in general position. Strictly speaking, we should only consider those times $`t`$ when $`S(t)`$ is in general position; when $`S(t)`$ is not in general position, we should decree that $`S(t)`$ is undefined, and have a discontinuity at $`t`$.

We shall see that we can go from $`S(0)`$ to $`S(T)`$ with only finitely many such discontinuous points. At such a discontinuity $`t`$, we still need to know how small we can make $`N_{S(t+\mathrm{\Delta }t)}N_{S(t\mathrm{\Delta }t)}`$ for small $`\mathrm{\Delta }t`$.

For small enough $`\mathrm{\Delta }t`$, the set $`S(t+\mathrm{\Delta }t)`$ is a very slight deformation of $`S(t)`$. What is missing is an understanding of what can make $`N_S`$ change as the set $`S`$ changes very slightly from $`S(t)`$ to $`S(t+\mathrm{\Delta }t)`$, and how small this change is. Let us answer this question.

Notice that, in the way we defined the deformation from $`S_{min}`$ to $`S_{max}`$, the points of $`S`$ moved only one at a time. Let us focus for now on the interval of time where $`P_1`$ moves towards $`Q_1`$.

Suppose that the number $`N_S`$ changes between time $`t`$ and time $`t+\mathrm{\Delta }t`$. Then it must be the case that for some $`i,j,k`$ and $`l`$ the circle $`P_iP_jP_k`$ contained (or did not contain) point $`P_l`$ at time $`t`$, but at time $`t+\mathrm{\Delta }t`$ it does not (or does) contain it. For this to be true, it must have happened that somewhere between times $`t`$ and $`t+\mathrm{\Delta }t`$, these four points must have been concyclic, or three of them must have been collinear. Since $`P_1`$ is the only point that has moved, we can conclude that $`P_1`$ must have crossed a circle or a line determined by the other points; this is what caused $`N_S`$ to change. We will call the circles and lines determined by the points $`P_2,P_3,\mathrm{},P_{2n+1}`$ the boundaries.

We can choose the path along which $`P_1`$ is going to move towards $`Q_1`$. To make things easier, we may assume that $`P_1`$ never crosses two of the boundaries at the same time. This can clearly be guaranteed: we know that these boundaries intersect pairwise in finitely many points, and all we have to do is avoid these intersection points in the path from $`P_1`$ to $`Q_1`$. We can also assume that $`\mathrm{\Delta }t`$ is small enough that $`P_1`$ crosses exactly one boundary between times $`t`$ and $`t+\mathrm{\Delta }t`$. Let us see how $`N_S`$ changes in this time interval.

It will be convenient to call a circle $`P_iP_jP_k`$ $`(a,b)`$-splitting (where $`a+b=2n2`$) if it has $`a`$ points of $`S`$ inside it and the remaining $`b`$ outside it. For example, an $`(n1,n1)`$-splitting circle is just a point-splitting circle.

First assume that $`P_1`$ crosses line $`P_iP_j`$, going from position $`P_1(t)=A`$ to position $`P_1(t+\mathrm{\Delta }t)=B`$. From the remarks made above, we know that only circle $`P_1P_iP_j`$ can change the value of $`N_S`$ by becoming or ceasing to be point- splitting. Assume that circle $`AP_iP_j`$ was $`(a,b)`$-splitting. Since $`P_1`$ only crossed the boundary $`P_iP_j`$ when going from $`A`$ to $`B`$, the region common to circles $`AP_iP_j`$ and $`BP_iP_j`$ cannot contain any points of $`S`$, as indicated in Figure 2. The region outside of both circles cannot contain points of $`S`$ either. For circle $`AP_iP_j`$ to be $`(a,b)`$-splitting, the other two regions must then contain $`a`$ and $`b`$ points respectively, as shown. Therefore, circle $`BP_iP_j`$ is $`(b,a)`$-splitting. It follows that $`AP_iP_j`$ was point-splitting if and only if $`BP_iP_j`$ is point-splitting (if and only if $`a=b=n1`$). We conclude that the value of $`N_S`$ doesn’t change when $`P_1`$ crosses a line determined by the other points; it can only change when $`P_1`$ crosses a circle.

Now assume that $`P_1`$ crosses circle $`P_iP_jP_k`$, going from position $`P_1(t)=A`$ inside the circle to position $`P_1(t+\mathrm{\Delta }t)=B`$ outside it. (The other case, when $`P_1`$ moves inside the circle, is analogous.) We can assume that $`P_1`$ crossed the arc $`P_iP_j`$ of the circle that doesn’t contain point $`P_k`$. Notice that $`A`$ must be outside triangle $`P_iP_jP_k`$ if we want $`P_1`$ to cross only one boundary in the time interval considered. Assume that circle $`AP_iP_j`$ was $`(a,b)`$-splitting. As before, we know that the only regions of Figure 3 containing points of $`S`$ are the one common to circles $`AP_iP_j`$ and $`BP_iP_j`$, and the one outside both of them. They must contain $`a1`$ and $`b`$ points respectively, for circle $`AP_iP_j`$ to be $`(a,b)`$-splitting. In this case, the value of $`N_S`$ can change only by circles $`P_iP_jP_k`$, $`P_1P_jP_k`$, $`P_1P_kP_i`$ and $`P_1P_iP_j`$ becoming or ceasing to be point-splitting. It is clear that circle $`P_iP_jP_k`$ went from being $`(a,b)`$-splitting to being $`(a1,b+1)`$-splitting. The same is true of circle $`P_1P_iP_j`$.

It is also not hard to see, by a similar argument, that circles $`P_1P_jP_k`$ and $`P_1P_kP_i`$ both went from being $`(a1,b+1)`$-splitting to being $`(a,b)`$-splitting. Again, the key assumption is that $`P_1`$ only crossed the boundary $`P_iP_jP_k`$ in this time interval.

So, by having $`P_1`$ cross circle $`P_iP_jP_k`$, we have traded two $`(a,b)`$-splitting and two $`(a1,b+1)`$-splitting circles for two $`(a1,b+1)`$-splitting and two $`(a,b)`$-splitting circles, respectively. It follows that the number $`N_S`$ of point-splitting circles remains constant when $`P_1`$ crosses a circle $`P_iP_jP_k`$ also.

We had shown that, as we moved $`P_1`$ to $`Q_1`$, $`N_S`$ could only possibly change in a time interval when $`P_1`$ crossed a boundary determined by the other points. But now we see that, even in such a time interval, $`N_S`$ does not change! Therefore moving $`P_1`$ to $`Q_1`$ doesn’t change the value of $`N_S`$. Similarly, moving $`P_i`$ to $`Q_i`$ doesn’t change $`N_S`$ either, for any $`1i2n+1`$. It follows that $`N_S`$ is the same for $`S_{min}`$ and $`S_{max}`$. In fact, $`N_S`$ is the same for any set $`S`$ of $`2n+1`$ points in general position!

## 3 $`N_S=n^2`$

Now that we know that the number $`N_S`$ depends only on the number of points in $`S`$, define $`N_{2n+1}`$ to be the number of point-splitting circles for a set of $`2n+1`$ points in general position. We compute $`N_{2n+1}`$ recursively.

Construct a set $`S`$ of $`2n+1`$ points as follows. First consider the vertices of a regular $`2n1`$-gon with center $`O`$. Now move them very slightly to positions $`P_1,\mathrm{},P_{2n1}`$ so that they are in general position. The difference will be so slight that all the lines $`OP_i`$ still split the remaining points into two sets of equal size, and all the circles $`P_iP_jP_k`$ still contain $`O`$. Also consider a point $`Q`$ which is so far away from the others that it lies outside of all the circles formed by the points considered so far. Of course, we need $`Q`$ to be in general position with respect to the remaining points. Let us count the number of point-splitting circles of $`S=\{O,P_1,\mathrm{},P_{2n1},Q\}`$.

First consider the circles of the form $`P_iP_jP_k`$. These circles contain $`O`$ and don’t contain $`Q`$; so they are point-splitting for $`S`$ if and only if they are point-splitting for $`\{P_1,\mathrm{},P_{2n1}\}`$. Thus there are $`N_{2n1}`$ such circles.

Next consider the circles $`OP_iP_j`$. It is clear that these circles contain at most $`n2`$ other $`P_k`$’s. They do not contain $`Q`$, so they contain at most $`n2`$ points, and they are not point-splitting.

Finally consider the circles that go through $`Q`$ and two other points $`X`$ and $`Y`$ of $`S`$. Circle $`QXY`$ splits the remaining points in the same way that line $`XY`$ does. More specifically, circle $`QXY`$ contains a point $`P`$ of $`S`$ if and only if $`P`$ is on the same side of line $`XY`$ that $`Q`$ is. This follows easily from the fact that $`Q`$ lies outside circle $`PXY`$. Therefore we have to determine how many lines determined by two of the points of $`S\{Q\}`$ split the remaining points of this set into two sets of $`n1`$ points each. This question is much easier; it is clear from our construction that the lines $`OP_i`$ do this and the lines $`P_iP_j`$ do not. It follows that the $`2n1`$ circles $`OP_iQ`$ are point-splitting, and the circles $`P_iP_jQ`$ are not.

Summarizing, the point-splitting circles of $`S`$ are the $`N_{2n1}`$ point-splitting circles of $`S\{O,Q\}`$ and the $`2n1`$ circles $`OP_iQ`$. Therefore $`N_{2n+1}=N_{2n1}+2n1`$. Since $`N_3=1`$, it follows inductively that $`N_{2n+1}=n^2.`$ This completes the proof of Theorem 1.3.

Theorem 3.1. Consider a set of $`2n+1`$ points in general position in the plane, and two non-negative integers $`a<b`$ such that $`a+b=2n2`$. There are exactly $`2(a+1)(b+1)`$ circles which are either $`(a,b)`$-splitting or $`(b,a)`$-splitting for the set of points.

Sketch of Proof. The argument of Section 2 carries directly to this situation, to show that the number of circles in consideration, which we denote $`N(a,b)`$, only depends on $`a`$ and $`b`$. Therefore it suffices to compute it recursively, using the set $`S`$ above. It is essential in the proof that $`a<n1`$.

Just as above, there are $`N(a1,b1)`$ such circles among the circles $`P_iP_jP_k`$. Among the $`OP_iP_j`$ there are exactly $`2n1`$ such circles, namely the circles $`OP_iP_{i+a+1}`$ (taking subscripts modulo $`2n1`$). There are also $`2n1`$ such circles among the $`QP_iP_j`$, namely the circles $`QP_iP_{i+a+1}`$. Finally, there are no such circles among the $`OP_iQ`$. Therefore $`N(a,b)=N(a1,b1)+4n2=N(a1,b1)+2a+2b+2`$.

Repeating the above argument for $`a=0`$, we get that $`N(0,b)=2b+2`$. If we combine this and the recursive relation obtained, Theorem 3.1 follows by induction.

It is worth mentioning at this point that Theorems 1.3 and 3.1 are closely related to a beautiful result of D.T. Lee, which gives a sharp bound for the number of vertices of an order $`j`$ Voronoi diagram. The connection is obtained if we embed our set $`S`$ of points on the surface of a sphere. Then the point-splitting circles of $`S`$ are put in correspondence with the “point-splitting planes” of a three-dimensional convex polytope with $`2n+1`$ vertices. These are known to be related to Voronoi diagrams. See \[1, p. 397\] for more details on this, and a proof of a result essentially equivalent to Theorems 1.3 and 3.1.

## 4 Questions

Our work completely determines the number of point-splitting circles, as well as the total number of $`(a,b)`$-splitting and $`(b,a)`$-splitting circles for a set of points in general position in the plane. However, we know very little about these numbers for sets of points that are not in general position.

The situation here is much more subtle. For example, the number of point-splitting circles of a set $`S`$ is not uniquely determined by the subsets of $`S`$ which are concyclic. Consider the following example. Let $`S_1`$ and $`S_2`$ be the two sets of seven points shown in Figure 4. Both of them are almost in general position; the only exception is that, for each of the two sets, there is a circle going through four points of the set. In $`S_1`$, this circle $`\mathrm{\Gamma }_1`$ contains exactly one point of $`S_1`$ inside it. In $`S_2`$, this circle $`\mathrm{\Gamma }_2`$ contains no points of $`S_2`$ inside it. In analogy with Theorem 1.3, where it did not matter which points were inside which circles, we might hope that $`S_1`$ and $`S_2`$ have the same number of point-splitting circles.

Unfortunately this is not the case. If we move $`A,B,C`$ or $`D`$ very slightly to put $`S_1`$ in general position, the resulting set will have $`9`$ point-splitting circles by Theorem 1.3. It is easy to see that exactly two of $`ABC,BCD,CDA`$ and $`DAB`$ are among these circles. When we deform the set back to $`S_1`$, these $`9`$ circles will still be point-splitting, but two of them will deform into $`\mathrm{\Gamma }_1`$. So $`S_1`$ has $`8`$ point-splitting circles.

Similarly, if we move $`a,b,c`$ or $`d`$ very slightly to put $`S_2`$ in general position, the resulting set will have $`9`$ point-splitting circles. But now we can see that when we deform the set back to $`S_2`$, none of these circles will deform into $`\mathrm{\Gamma }_2`$, because $`\mathrm{\Gamma }_2`$ contains no points of $`S_2`$. Therefore $`S_2`$ has $`9`$ point-splitting circles.

Even if the number of point-splitting circles is not constant, we might be able to say something about it. As a small example, consider all sets of seven points which are almost in general position, except that four of them are concyclic. It is possible to show, by an argument similar to the above, that such a set can only have $`8`$ or $`9`$ point-splitting circles. It seems reasonable that, in general, one might be able to define some measure of how far a set $`S`$ is from being in general position, and to obtain bounds for $`N_S`$ in terms of that measure.

## 5 Acknowledgements

The results in this article were first presented by the author at Dan Kleitman’s 65th birthday conference at the Massachusetts Institute of Technology in 1999. The author wishes to thank Dan Kleitman and Richard Stanley for useful discussions on this subject during the conference.

He also wishes to thank Timothy Chow, who took an interest in this result and mentioned it to several people. It is due to his efforts that the author became aware of the equivalent result in  and the connection with Voronoi diagrams.