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# Cosmological solutions in multidimensional model with multiple exponential potential

## 1 Introduction

Recently, the discovery of the cosmic acceleration  stimulated a lot of papers on multidimensional cosmology with the aim to explain this phenomenon using certain multidimensional models, e.g. those of superstring or supergravity origin (see -, and references therein). These solutions usually deal with time-dependent scale factors of internal spaces, thus overcoming a ”no-go” theorem for static (and compact) internal spaces . It should be noted that certain part of publications does not deal with new exact solutions but use old ones (sometimes rediscovered or written in different parametrization).

A simple example of vacuum solution with acceleration was considered by Townsend and Wohlfarth in . This is a $`(4+n)`$-dimensional solution to vacuum Einstein equations with $`n`$-dimensional internal space of negative curvature and $`4`$-dimensional ”our space” containing expanding $`3`$-dimensional flat subspace. The solution from  is a special solution, in a special model with a very special choice of conformal frame. We note that it is a special case of more general vacuum solution from , describing the ”evolution” of $`(n1)`$ Ricci-flat spaces and one Einstein space of non-zero curvature. The solution from  was generalized to scalar-vacuum case in  and to the case of composite p-brane configurations in  (for non-composite case see also ). For a review see also . We note that the solutions from  may be also used for generation of special solutions with several curved factor spaces (e.g. from ) using a curvature-splitting trick  (when Einstein space of non-zero curvature is chosen as a product of several Einstein spaces). (For three integrable classes of vacuum cosmological solutions with two factor spaces of non-zero curvature in dimensions $`D=10,11`$ see .)

At present rather popular models are those with multiple exponential potential of the scalar fields (see, for example,  and refs. therein).

Such potentials also arise naturally in certain supergravitational models  and in sigma-models , related to configurations with $`p`$-branes.

Here we consider the $`D`$-dimensional model governed by the action

$`S_{act}={\displaystyle _M}d^Dz\sqrt{|g|}\{[g]h_{\alpha \beta }g^{MN}_M\phi ^\alpha _N\phi ^\beta 2V_\phi (\phi )\}+S_{GH},`$ (1.1)
with the scalar potential

$$V_\phi (\phi )=\underset{sS}{}\mathrm{\Lambda }_s\mathrm{exp}[2\lambda _s(\phi )].$$
(1.2)
Here $`S_{\mathrm{GH}}`$ is the standard Gibbons-Hawking boundary term .

The notations used here are the following ones:

* $`\phi =(\phi ^\alpha )`$ is the vector from scalar fields in the space $`\text{IR}^l`$ with a metric determined by non-degenerate $`l\times l`$ matrix $`(h_{\alpha \beta })`$ with inverse one $`(h^{\alpha \beta })`$,
$`\alpha =1,\mathrm{},l`$;
* $`\mathrm{\Lambda }_s`$ are constant terms,
$`sS`$;
* $`\lambda _s`$ is an $`1`$-form on $`\text{IR}^l`$: $`\lambda _s(\phi )=\lambda _{s\alpha }\phi ^\alpha `$$`\lambda _s^\alpha =h^{\alpha \beta }\lambda _{s\beta }`$;
* $`g=g_{MN}dz^Mdz^N`$ is the metric, $`|g|=|det(g_{MN})|`$,
$`M,N`$ are world indicies
(that may be numerated by $`1,\mathrm{},D`$);
* $`i,j=0,\mathrm{},n`$ are indicies describing a chain of factor spaces;
$`A=i,\alpha `$ and $`B=j,\beta `$ are minisuperspace indicies
(that may be numerated also by $`0,\mathrm{},n+l`$).

In this article we obtain new (and general) families of classical and quantum “cosmological” solutions with vector coupling constants obeying

$$\lambda _{s\alpha }\lambda _{s^{}\beta }h^{\alpha \beta }=\frac{D1}{D2}$$
(1.3)
for $`ss^{}`$ and

$$\lambda _{s\alpha }\lambda _{s\beta }h^{\alpha \beta }\frac{D1}{D2}.$$
(1.4)
Here we find exact solutions with scale factors and scalar fields depending upon one variable $`u`$ (”time”). We keep the parameter $`w=\pm 1`$ in the metric $`g=w\text{e}^{2\gamma (u)}dudu+\mathrm{}`$ for a future consideration of static configurations . They may be of interest in a context of ”black hole” solutions with unusual asymptotics (see  and references therein). We single out two new subclasses of solutions with power-law and exponential behaviour of scale factors, e.g. those with accelerated expansion. A new result here is generalized isotropization behaviours of certain classes of classical solutions that give new examples of asymptotically accelerated expansion.

The paper is organized as follows. In Section 2 the model is formulated. In Section 3 a class of classical cosmological solutions corresponding to orthogonal $`U`$-vectors is presented. Section 4 deals with special solutions exhibiting power-law and exponential dependence of scale factors. Here we describe the solutions with accelerated expansion. A special subsection is devoted to a generalized isotropization. In section 5 quantum analogues of the classical solutions are considered. In Appendix the equations of motion, corresponding to (1.1) are written (Appendix A) and derivations of special solutions of power law and exponentional type are presented (Appendix B and C); also a classification of vectors in Euclidean space obeying relations (1.3) and (1.4) is done (Appendix D).

We note that the classical and quantum solutions for cosmological constant case (i.e. for one term in potential with $`\lambda _s=0`$) were considered earlier in  and  for vacuum and scalar-vacuum cases, respectively.

## 2 The model

Let

$$M=\text{IR}_.\times M_0\times \mathrm{}\times M_n$$
(2.1)
be a manifold equipped with the metric

$$g=w\text{e}^{2\gamma (u)}dudu+\underset{i=0}{\overset{n}{}}\text{e}^{2\varphi ^i(u)}g^i,$$
(2.2)
where $`w=\pm 1`$, $`\text{IR}_.`$ is open interval in IR, $`u`$ is a distinguished coordinate; $`g^i`$ is a Ricci-flat metric on $`d_i`$-dimensional manifold $`M_i`$:

$$_{m_in_i}[g^i]=0,$$
(2.3)
$`d_i=dimM_i`$, $`i=0,\mathrm{},n`$; $`n𝐍`$.

For dilatonic scalar fields we put

$$\phi ^\alpha =\phi ^\alpha (u),$$
(2.4)
It may be verified that the equations of motion corresponding to (1.1) for the field configuration (2.2), (2.4) are equivalent to equations of motion for 1-dimensional $`\sigma `$-model with the action

$$S_\sigma =\frac{1}{2}𝑑u𝒩\left\{G_{ij}\dot{\varphi }^i\dot{\varphi }^j+h_{\alpha \beta }\dot{\phi }^\alpha \dot{\phi }^\beta 2𝒩^2V\right\},$$
(2.5)
where $`\dot{x}dx/du`$,

$$V=wV_\phi (\phi )\text{e}^{2\gamma _0(\varphi )}$$
(2.6)
is the potential ($`V_\phi `$ is defined in (1.2)) with

$$\gamma _0(\varphi )\underset{i=0}{\overset{n}{}}d_i\varphi ^i,$$
(2.7)
and

$$𝒩=\mathrm{exp}(\gamma _0(\varphi )\gamma )>0$$
(2.8)
is lapse function. Here

$$G_{ij}=d_i\delta _{ij}d_id_j,G^{ij}=\frac{\delta ^{ij}}{d_i}+\frac{1}{2D},$$
(2.9)
$`i,j=0,\mathrm{},n`$, are components of a gravitational part of minisupermetric and its dual .

### 2.1 Minisuperspace notations

In what follows we consider minisuperspace $`\text{IR}^{n+1+l}`$ with points

$$x(x^A)=(\varphi ^i,\phi ^\alpha )$$
(2.10)
equipped by minisuperspace metric $`\overline{G}=\overline{G}_{AB}dx^Adx^B`$ defined by the matrix and inverse one as follows:

$`(\overline{G}_{AB})=\left(\begin{array}{cc}G_{ij}& 0\\ 0& h_{\alpha \beta }\end{array}\right),(\overline{G}^{AB})=\left(\begin{array}{cc}G^{ij}& 0\\ 0& h^{\alpha \beta }\end{array}\right).`$ (2.15)
The potential (2.6) reads

$$V=w\underset{sS}{}\mathrm{\Lambda }_s\text{e}^{2U^s(x)}$$
(2.16)
where $`U^s(x)=U_A^sx^A`$ are defined as follows

$$U^s=U^s(\varphi ,\phi )=\lambda _{s\alpha }\phi ^\alpha +\underset{i=0}{\overset{n}{}}d_i\varphi ^i,$$
(2.17)
or in components

$$(U_A^s)=(d_i,\lambda _{s\alpha }).$$
(2.18)
The integrability of the Lagrange system (2.5) depends upon the scalar products of co-vectors $`U^s`$ corresponding to $`\overline{G}`$:

$$(U,U^{})=\overline{G}^{AB}U_AU_B^{},$$
(2.19)
These products have the following form

$$(U^s,U^s^{})=b+\lambda _s\lambda _s^{},$$
(2.20)
where $`s,s^{}S`$ and

$`\lambda _s\lambda _s^{}\lambda _{s\alpha }\lambda _{s^{}\beta }h^{\alpha \beta },`$ (2.21)
$`b={\displaystyle \frac{D1}{D2}}.`$ (2.22)
We put the orthogonality restriction on the vectors $`U^s`$

$$(U^s,U^s^{})=b+\lambda _s\lambda _s^{}=0,$$
(2.23)
$`ss^{}`$, and also consider non-degenerate case

$$(U^s,U^s)=b+\lambda _s\lambda _s0$$
(2.24)
(see (1.3) and (1.4)).

In what follows we denote

$`h_s=(U^s,U^s)^1{\displaystyle \frac{1}{\lambda _s^2b}}.`$ (2.25)
The further consideration is based upon the orthogonality conditions assumed.

## 3 Classical solutions

Here we will integrate the Lagrange equations corresponding to the action (2.5) in the harmonic time gauge $`\gamma =\gamma _0`$. We get a Lagrangian

$$L=\frac{1}{2}\overline{G}_{AB}\dot{x}^A\dot{x}^BV,$$
(3.1)
where $`(\overline{G}_{AB})`$ and $`V`$ are defined in (2.15) and (2.16), respectively. The zero-energy constraint reads:

$$E=\frac{1}{2}\overline{G}_{AB}\dot{x}^A\dot{x}^B+V=0.$$
(3.2)
The solutions to Lagrange equations corresponding to (3.1) in the orthogonal and non-degenerate case (2.23), (2.24) read (see ):

$$x^A=\underset{sS}{}\frac{U^{sA}}{(U^s,U^s)}\mathrm{ln}|f_s|+c^Au+\overline{c}^A,$$
(3.3)
Here

$`f_s=R_s\mathrm{sh}(\sqrt{\mathrm{C}_\mathrm{s}}(\mathrm{u}\mathrm{u}_\mathrm{s})),\mathrm{C}_\mathrm{s}>0,ϵ_\mathrm{s}<0;`$ (3.4)
$`R_s\mathrm{sin}(\sqrt{|C_s|}(uu_s)),C_s<0,ϵ_s<0;`$ (3.5)
$`R_s\mathrm{ch}(\sqrt{\mathrm{C}_\mathrm{s}}(\mathrm{u}\mathrm{u}_\mathrm{s})),\mathrm{C}_\mathrm{s}>0,ϵ_\mathrm{s}>0;`$ (3.6)
$`|2\mathrm{\Lambda }_s/h_s|^{1/2}(uu_s),C_s=0,ϵ_s<0,`$ (3.7)
where $`u_s`$ and $`C_s`$ are constants and

$`R_s=|2\mathrm{\Lambda }_s/(h_sC_s)|^{1/2},`$ (3.8)
$`ϵ_s=\text{sign}(w\mathrm{\Lambda }_sh_s).`$ (3.9)
Vectors $`c=(c^A)`$ and $`\overline{c}=(\overline{c}^A)`$ satisfy the linear constraint relations

$`U^s(c)=U_A^sc^A={\displaystyle \underset{i=0}{\overset{n}{}}}d_ic^i+\lambda _{s\alpha }c^\alpha =0,`$ (3.10)
$`U^s(\overline{c})=U_A^s\overline{c}^A={\displaystyle \underset{i=0}{\overset{n}{}}}d_i\overline{c}^i+\lambda _{s\alpha }\overline{c}^\alpha =0.`$ (3.11)
The zero-energy constraint reads

$$E=\underset{sS}{}E_s+\frac{1}{2}\overline{G}_{AB}c^Ac^B=0,$$
(3.12)
where $`C_s=2E_s(U^s,U^s)`$, or, equivalently,

$$\underset{sS}{}C_sh_s+h_{\alpha \beta }c^\alpha c^\beta +\underset{i=0}{\overset{n}{}}d_i(c^i)^2\left(\underset{i=0}{\overset{n}{}}d_ic^i\right)^2=0.$$
(3.13)
The solutions.

Using (3.3) and relations for contravariant components $`U^{sA}=\overline{G}^{AB}U_B^s`$

$$U^{si}=G^{ij}U_j^s=\frac{1}{D2},U^{s\alpha }=\lambda _s^\alpha ,$$
(3.14)
we are led to relations for the metric

$`g=\left({\displaystyle \underset{sS}{}}f_s^{2bh_s}\right)\left\{\text{e}^{2cu+2\overline{c}}wdudu+\left({\displaystyle \underset{sS}{}}f_s^{2h_s}\right){\displaystyle \underset{i=0}{}}\text{e}^{2c^iu+2\overline{c}^i}g^i\right\}.`$ (3.15)
and scalar fields

$$\phi ^\alpha =\underset{sS}{}h_s\lambda _s^\alpha \mathrm{ln}|f_s|+c^\alpha u+\overline{c}^\alpha ,$$
(3.16)
$`\alpha =1,\mathrm{},l`$.

Here

$$c=\underset{i=0}{\overset{n}{}}d_ic^i,\overline{c}=\underset{i=0}{\overset{n}{}}d_i\overline{c}^i.$$
(3.17)
The functions $`f_s`$ are defined in (3.4)–(3.7) and the relations on the parameters of solutions $`c^A`$, $`\overline{c}^A`$ $`(A=i,\alpha )`$ and $`C_s`$ imposed in (3.10)–(3.11) and (3.13), respectively. ($`h_s`$ are defined in (2.25).)

## 4 Special solutions

Now we consider a special case of classical solutions from the previous section when $`C_s=u_s=c^i=c^\alpha =0`$ and

$$w\mathrm{\Lambda }_s(\lambda _s^2b)>0,$$
(4.1)
$`sS`$.

We get two families of solutions written in synchronous-type variable with:

A) power-law dependence of scale factors for $`B1`$,

B) exponential dependence of scale factors for $`B=1`$,

where

$$B=B(\lambda )=b\underset{sS}{}h_s.$$
(4.2)
Note that $`h_s=h_s(\lambda _s)=(\lambda _s^2b)^1`$.

### 4.1 Power-law solutions

Let us consider the solution corresponding to the case $`B1`$. The solution reads (see Appendix B)

$`g=wd\tau d\tau +\tau ^{2\nu }{\displaystyle \underset{i=0}{\overset{n}{}}}A_ig^i,`$ (4.3)
$`\phi ^\alpha ={\displaystyle \frac{1}{B+1}}{\displaystyle \underset{sS}{}}h_s\lambda _s^\alpha \mathrm{ln}\tau +\phi _0^\alpha ,`$ (4.4)
where $`\tau >0`$,

$$\nu =\frac{B}{(B+1)(D1)},$$
(4.5)
and

$$|2\mathrm{\Lambda }_s/h_s|\mathrm{exp}(2\lambda _{s\alpha }\phi _0^\alpha )=(B+1)^2,$$
(4.6)
$`sS`$; and $`A_i>0`$ are arbitrary constants.

Solutions with ”acceleration”. Let us consider the cosmological case $`w=1`$. We get an accelerated expansion of factor spaces if and only if $`\nu >1`$ or, equivalently,

$$b<B(\lambda )<1.$$
(4.7)
Remark 1.  For $`1<B=B(\lambda )<0`$ we get $`\nu <0`$, that corresponds to a contraction of factor spaces. When $`B<b`$ or $`B>0`$ we are led to the expansion with deceleration, since $`0<\nu <1`$ in this case. The case $`B=0`$ corresponds to static toy “universe”. For $`B=b`$ we obtain the expansion with zero acceleration ($`\nu =1`$).

Let us consider the case when the matrix $`(h_{\alpha \beta })`$ is positive definite (e.g. $`h_{\alpha \beta }=\delta _{\alpha \beta }`$). It was shown in Appendix D that the condition $`B(\lambda )1`$ in this case implies that vectors $`\lambda _s`$ are linearly independent. It follows also from Appendix D that the condition $`B(\lambda )<1`$ implies that

$$\lambda _{s_0}^2<b,\lambda _s^2>b\mathrm{for}\mathrm{all}ss_0.$$
(4.8)
Here $`s_0S`$. These relations combined with inequalities (4.1) lead to

$$\mathrm{\Lambda }_{s_0}>0,\mathrm{\Lambda }_s<0\mathrm{for}\mathrm{all}ss_0,$$
(4.9)
i.e. one exponent term in potential should be positive and others should be negative.

### 4.2 Solutions with exponential scale factors

Here we consider the solution corresponding to the case $`B=1`$. The solution reads (see Appendix C)

$`g=wd\tau d\tau +\mathrm{exp}(2m\tau ){\displaystyle \underset{i=0}{\overset{n}{}}}A_ig^i,`$ (4.10)
$`\phi ^\alpha =(D1)m\tau {\displaystyle \underset{sS}{}}h_s\lambda _s^\alpha +\phi _0^\alpha ,`$ (4.11)
where

$$|2\mathrm{\Lambda }_s/h_s|\mathrm{exp}(2\lambda _{s\alpha }\phi _0^\alpha )=m^2(D1)^2.$$
(4.12)
$`sS`$; $`m`$ is parameter and $`A_i>0`$ are arbitrary constants.

In the cosmological case $`w=1`$ we get an accelerated expansion of factor spaces if and only if $`m>0`$.

Let us consider the case when $`h_{\alpha \beta }`$ is positive definite. As it was shown in Appendix D the condition $`B(\lambda )=1`$ in this case implies that vectors $`\lambda _s`$ are linearly dependent, the linear term in (4.11) vanishes and hence

$$\phi ^\alpha =\phi _0^\alpha ,$$
(4.13)
i.e. for $`(h_{\alpha \beta })`$ of Euclidean signature all scalar fields are constants when exponential expansion is under consideration. We note that any subset of vectors $`\lambda _s,sS\{s_0\}`$ is linearly independent one in this case (see Appendix D).

Remark 2. It may be shown that for the solution (4.10)-(4.12) the kinetic term for scalar fields vanishes, i.e. $`h_{\alpha \beta }g^{MN}_M\phi ^\alpha _N\phi ^\beta =0`$ and the scalar potential is constant: $`V_\phi (\phi )=V_\phi (\phi _0)\mathrm{\Lambda }`$ and $`(w)2\mathrm{\Lambda }=m^2(D1)(D2)`$ (see (4.1)). For $`w=1`$ and $`\mathrm{\Lambda }>0`$ the metric (4.10) is coinciding with that for multidimensional model from  with pure cosmological term $`\mathrm{\Lambda }`$. Thus, the exponential expansion of factor spaces is driven by effective cosmological term $`\mathrm{\Lambda }`$ and is not sensitive to time dependence of scalar fields, that takes place only for non-Euclidean signatures of $`(h_{\alpha \beta })`$.

### 4.3 One-exponent example

The simplest case with one term in the potential (1.2) is interesting due to the absence of the restriction (2.23). It follows from eq. (4.1) that the potential is positive:

$$V_\phi =\mathrm{\Lambda }\text{e}^{2\lambda _\alpha \phi ^\alpha },\mathrm{\Lambda }>0$$
(4.14)
(we discard the indices in this example).

The accelerated power law expansion takes place, if

$$0<\lambda ^2<\frac{1}{D2},$$
(4.15)
i.e. $`\lambda ^2`$ should be small enough but not equal to zero. The inequality is coinciding (up to notations) with the eq. (1.9) from . In this case (4.5) reads as $`\nu =(\lambda ^2(D2))^1`$.

For the exponential expansion eqs. (4.1) and $`B(\lambda )=1`$ imply

$$\lambda ^2=0,$$
(4.16)
i.e. $`\lambda `$ should be zero vector for positive definite matrix $`(h_{\alpha \beta })`$ or light-like vector for $`(h_{\alpha \beta })`$ of pseudo-Euclidean signature.

### 4.4 Generalized isotropization

In the case of one exponent the special solutions considered above are attractors in the limit $`\tau +\mathrm{}`$ ($`\tau `$ is ”synchronous” variable) for general solutions with sinh-dependent function $`f_s`$ (see (3.4)), when $`ϵ_s<0`$, or, equivalently, (4.1) is satisfied and $`B1`$. This fact may be readily verified using the relation $`f_s(uu_s)`$ for $`uu_s`$.

Here we consider a more general case of this attractor behavior for multi-exponent case. Let $`S_{}`$ be a subset of all $`S`$ satisfying $`ϵ_s<0`$, or, equivalently, relations (4.1). Let us study a family of solutions with positive $`C_s`$ and hence with $`\mathrm{sinh}`$-dependence of $`f_s`$ for $`sS_{}`$ and $`\mathrm{cosh}`$-dependence of $`f_s`$ for $`sSS_{}`$. Consider the solutions on the interval $`(u_{},+\mathrm{})`$, where

$$u_{}=\mathrm{max}(u_s,sS_{}).$$
(4.17)
We denote by $`S_{}`$ a subset of all $`sS_{}`$ obeying $`u_s=u_{}`$. All functions $`f_s(uu_s)`$ are smooth on this interval and near $`u_{}`$ they behave as $`f_s(uu_{})`$ for $`sS_{}`$, or as (non-zero) constants for $`sSS_{}`$.

Introducing the synchronous time variable by relation

$$\tau =𝑑u\mathrm{exp}(cu+\overline{c})\underset{s}{}(f_s(uu_s))^{bh_s}+const$$
(4.18)
we get $`\tau +\mathrm{}`$ as $`uu_{}`$ in two cases:

A) $`B_{}<1`$ and

$$\tau (uu_{})^{B_{}+1};$$
(4.19)
B) $`B_{}=1`$ and

$$\tau \mathrm{ln}|uu_{}|.$$
(4.20)
Here we denote

$$B_{}=B(\lambda ,S_{})=\underset{sS_{}}{}h_sb.$$
(4.21)
Case A. Using arguments analogous to those presented in Appendix B we get the following asymptotical relations in the limit $`\tau +\mathrm{}`$ for $`B_{}<1`$

$`g_{as}=wd\tau d\tau +{\displaystyle \underset{i=0}{\overset{n}{}}}A_i\tau ^{2\nu _{}}g^i,`$ (4.22)
$`\phi _{as}^\alpha ={\displaystyle \frac{1}{B_{}+1}}{\displaystyle \underset{sS_{}}{}}h_s\lambda _s^\alpha \mathrm{ln}\tau +\phi _0^\alpha ,`$ (4.23)
where

$$\nu _{}=\frac{B_{}}{(B_{}+1)(D1)},$$
(4.24)
and $`A_i>0`$ are constants.

Thus, we get an isotropization behaviour of scale factors for $`\tau +\mathrm{}`$, that is supported by exponential terms in the potential labelled by $`sS_{}`$. Other exponential terms do not contribute to the power-law index $`\nu _{}`$.

For

$$b<B_{}<1$$
(4.25)
we obtain the accelerated expansion for large enough values of $`\tau `$.

Case B.

An analogous consideration in the case $`B_{}=1`$ leads to the following asymptotical relations in the limit $`\tau +\mathrm{}`$ (see also Appendix C)

$`g_{as}=wd\tau d\tau +{\displaystyle \underset{i=0}{\overset{n}{}}}A_i\mathrm{exp}(2m\tau )g^i,`$ (4.26)
$`\phi _{as}^\alpha =(D1)m\tau {\displaystyle \underset{sS_{}}{}}h_s\lambda _s^\alpha +\phi _0^\alpha ,`$ (4.27)
where $`m>0`$ and

$$|2\mathrm{\Lambda }_s/h_s|\mathrm{exp}(2\lambda _{s\alpha }\phi _0^\alpha )=m^2(D1)^2.$$
(4.28)
$`sS_{}`$. Thus, we obtain the accelerated expansion for large enough values of $`\tau `$.

This is another regime of isotropization behaviour (for one-exponent case with $`\lambda _s=0`$ see also ).

Remark 3. It may be shown that for the asymptotical solution (4.26)-(4.28) the kinetic term for scalar fields asymptotically vanishes, i.e. $`h_{\alpha \beta }g^{MN}_M\phi ^\alpha _N\phi ^\beta 0`$ as $`\tau +\mathrm{}`$ and the scalar potential is asymptotically constant:

$$V_\phi (\phi )\underset{sS_{}}{}\mathrm{\Lambda }_s\mathrm{exp}[2\lambda _s(\phi _0)]\mathrm{\Lambda }_{}$$
(4.29)
as $`\tau +\mathrm{}`$. Here $`(w)2\mathrm{\Lambda }_{}=m^2(D1)(D2)`$. For $`w=1`$ ans $`\mathrm{\Lambda }_{}>0`$ the “asymptotical” metric (4.10) is coinciding with that for multidimensional model from  with pure cosmological term $`\mathrm{\Lambda }_{}`$. When $`(h_{\alpha \beta })`$ has Euclidean signature the exponential isotropization behaviour takes place only for $`S=S_{}`$ (it follows from Appendix D). In this case $`\phi _{as}=\phi _0`$.

### 4.5 Kasner-type behaviour

Let us study the solution with $`\mathrm{sinh}`$-dependence in the limit $`u\mathrm{}`$. We restrict ourselves by Kasner-like behaviour of solutions for small $`\tau `$. The functions (3.4) and (4.18) behave as

$$f_s\text{e}^{\sqrt{C_s}u},\tau \text{e}^{(b\mathrm{\Sigma }+c)u},u+\mathrm{},$$
(4.30)
where $`\mathrm{\Sigma }=\underset{sS}{}h_s\sqrt{C_s}`$. To get the limit $`\tau +0`$ we put

$$b\mathrm{\Sigma }+c<0.$$
(4.31)
The metric and scalar fields in the limit $`\tau +0`$ read as follows:

$`g_{as}=wd\tau d\tau +{\displaystyle \underset{i=0}{\overset{n}{}}}\stackrel{~}{A}_i\tau ^{2\alpha ^i}g^i,`$ (4.32)
$`\phi _{as}^\beta =\alpha ^\beta \mathrm{ln}\tau +\stackrel{~}{\phi }_0^\beta ,`$ (4.33)
where the Kasner parameters are given by the formulas

$`\alpha ^i={\displaystyle \frac{\mathrm{\Sigma }(D2)^1+c^i}{b\mathrm{\Sigma }+c}},`$ (4.34)
$`\alpha ^\beta ={\displaystyle \frac{\underset{sS}{}h_s\sqrt{C_s}\lambda _s^\beta +c^\beta }{b\mathrm{\Sigma }+c}}.`$ (4.35)
One can see that these parameters obey the generalized Kasner relations:

$`{\displaystyle \underset{i=0}{\overset{n}{}}}d_i\alpha ^i=1,{\displaystyle \underset{i=0}{\overset{n}{}}}d_i(\alpha ^i)^2+h_{\alpha \beta }\alpha ^\alpha \alpha ^\beta =1.`$ (4.36)
We note that for $`b\mathrm{\Sigma }+c<0`$ we get a Kasner-like asymptotics in the limit $`\tau \mathrm{}`$ (with $`\tau `$ replaced by $`|\tau |`$ in (4.32) and (4.33)).

#### 4.5.1 Example: two spaces with one scalar field

As an example let us consider the model with two factor spaces, one scalar field and one term in potential. It is not difficult to verify, that the parameters

$`\alpha ^0={\displaystyle \frac{(d_0R)}{d_0(d_0+d_1)}},\alpha ^1={\displaystyle \frac{(d_1+R)}{d_1(d_0+d_1)}},\alpha ^\beta =0,`$ (4.37)
where $`R=\sqrt{d_0d_1(D2)}`$, satisfy the eqs. (4.36).

For $`d_1=3`$, $`d_2=6`$ the behaviour of scale factors are presented in Fig. 1, when: a) $`\lambda ^2=1/16`$ (asymptotically power law expansion); b) $`\lambda =0`$ (asymptotically exponential expansion).

Remark 4. It will be shown in a separate publication that Kasner-like behaviour near the singularity (as $`\tau +0`$) is a generic one for $`n>0`$ and a wide class of scalar potentials (1.2).

## 5 Quantum solutions

Here we find the quantum analogues of the classical solutions presented above.

### 5.1 General solutions

The standard quantization of the energy constraint $`\mu E=0`$, where $`\mu 0`$ is a parameter, leads to the Wheeler-DeWitt (WDW) equation (see, for example, ) in the harmonic-time gauge $`\gamma =\gamma _0(\varphi )`$

$$\widehat{H}\mathrm{\Psi }\left(\frac{1}{2\mu }\mathrm{\Delta }[\overline{G}]+\mu V\right)\mathrm{\Psi }=0.$$
(5.1)
Here $`V`$ is potential defined in (2.16), $`\mathrm{\Delta }[\overline{G}]`$ is the Laplace operator, corresponding to $`\overline{G}`$, and $`\mathrm{\Psi }`$ is wave function.

The minisuperspace metric may be diagonalized by the linear transformation

$$z^{\overline{A}}=S^{\overline{A}}{}_{B}{}^{}x_{}^{B},(z^{\overline{A}})=(z^0,z^a,z^s)$$
(5.2)
as follows

$$\overline{G}_{AB}dx^Adx^B=dz^0dz^0+\underset{sS}{}\eta _sdz^sdz^s+dz^adz^b\eta _{ab},$$
(5.3)
where $`a,b=1,\mathrm{},n+l|S|`$; $`\eta _{ab}=\eta _{aa}\delta _{ab}`$; $`\eta _{aa}=\pm 1`$, $`\eta _s=\text{sign}(h_s)`$, and

$$q_sz^s=U^s(x),$$
(5.4)
with

$$q_s=|h_s|^{1/2}=\sqrt{|\lambda _s^2b|}>0,$$
(5.5)
$`sS`$.

We are seeking the solution to WDW equation (5.1) by the method of the separation of variables, i.e. we put

$$\mathrm{\Psi }=\left(\underset{sS}{}\mathrm{\Psi }_s(z^s)\right)\text{e}^{\mathrm{i}p_az^a}.$$
(5.6)
The wave function (5.6) satisfies to WDW equation (5.1) if

$$2\widehat{H}_s\mathrm{\Psi }_s\left\{\eta _s\frac{}{z^s}\left(\frac{}{z^s}\right)2w\mathrm{\Lambda }_s\text{e}^{2q_sz^s}\right\}\mathrm{\Psi }_s=2_s\mathrm{\Psi }_s,$$
(5.7)
where

$$\eta ^{ab}p_ap_b+2\underset{sS}{}_s=0.$$
(5.8)
(compare with classical relation (3.12)).

The linearly independent solutions to eq. (5.7) reads

$$\mathrm{\Psi }_s(z^s)=B_{\omega _s}^s\left(\sqrt{2w\eta _s\mathrm{\Lambda }_s}\frac{\text{e}^{q_sz^s}}{q_s}\right),$$
(5.9)
where

$$\omega _s=\sqrt{2_sh_s},$$
(5.10)
$`sS`$ and $`B_\omega ^s=a_sI_\omega +b_sK_\omega `$ are superpositions of modified Bessel function.

The general solution to the WDW equation (5.1) is a superposition of the ”separated” solutions (5.6).

### 5.2 Special solutions

Let us consider the special ”ground state” solutions with $`p_a=_s=0`$ in (5.9), when relations (4.1) are imposed. These solutions read

$$\mathrm{\Psi }=\underset{sS}{}\left(a_sJ_0(\sqrt{2|h_s\mathrm{\Lambda }_s|}v\text{e}^{\lambda _s(\phi )})+b_sH_0^{(1)}(\sqrt{2|h_s\mathrm{\Lambda }_s|}v\text{e}^{\lambda _s(\phi )})\right),$$
(5.11)
where $`J_0`$ and $`H_0^{(1)}`$ are the Bessel and Hankel functions, respectively, and $`v=\mathrm{exp}(\underset{i=0}{\overset{n}{}}d_i\varphi ^i)`$ is a volume scale factor. These quantum solutions correspond to special classical solutions considered in the previous section.

We note that for large values of

$$v_s=\sqrt{2|h_s\mathrm{\Lambda }_s|}v\mathrm{exp}(\lambda _s(\phi ))$$
(5.12)
we get a quasi-classical regime, that may be obtained (along a line as it was done in ) using asymptotical relations for Bessel functions. For small values of ”quasivolumes” $`v_s`$ we are led to a quantum regime. The crucial point here is that the quantum domain is not defined only by volume scale factor $`v`$, but by ”quasivolumes” $`v_s`$, depending also upon scalar fields. In strong enough scalar fields with certain signs (and the direction of vector $`\phi `$) and/or small $`|h_s\mathrm{\Lambda }_s|`$ one can obtain a quantum behaviour of a wave function for big enough values of volume scale factor $`v`$. Analogous ”effect” takes place also for quantum solutions with ”perfect fluid”  when certain equations of state are adopted.

## 6 Conclusions

Here we obtained a family of multidimensional cosmological solutions with $`(n+1)`$ Ricci-flat spaces in the theory with several scalar fields and multiple exponential potential.

The classical and quantum solutions are obtained if the (orthogonality and non-degeneracy) relations (2.23) and (2.24) on $`U`$-vectors are imposed, or, equivalently, when relations (1.3) and (1.4) on coupling vectors are satisfied. These solutions in fact correspond to $`A_1+\mathrm{}+A_1`$ Toda-like solutions (for perfect fluid case see ).

Here we singled out the solutions with power-law and exponential behaviours of scale factors. We proved that power-law dependence may take place only when coupling vectors $`\lambda _s`$ are linearly independent; exponential dependence may occur only for a linearly dependent set of $`\lambda _s`$ obeying the condition $`B(\lambda )=1`$. We obtained the restriction on coupling vectors (4.7) that cuts power-law solutions with acceleration. In subsection 4.3 the generalized isotropization behavior of certain class of general solutions was found. Any asymptotics with power-law or exponential behavior corresponds to a certain subset of exponential terms in the potential (all other exponential terms do not contribute to this asymptotical behaviour).

In the quantum case we solved the Wheeler-DeWitt equation and singled out special ”ground state” solutions that are quantum analogous of special classical solutions from Section 4. These solutions are defined as products of Bessel functions and depend upon a set of quasivolumes that may be small enough (i.e. belonging to ”quantum domain”) even when the volume scale factor of the toy ”universe” is in the classical region (in the Planck scale).

###### Acknowledgments.

This work was supported in part by the Russian Ministry of Science and Technology, Russian Foundation for Basic Research (RFFI-01-02-17312-a) and Project DFG (436 RUS 113/678/0-1(R)).

## Appendix

### A Equations of motion

Here we outline for the sake of completeness the equations of motions corresponding to the action (1.1)

$`_{MN}{\displaystyle \frac{1}{2}}g_{MN}=T_{MN},`$ (A.1)
$`[g]\phi ^\alpha {\displaystyle \underset{sS}{}}2\lambda _s^\alpha e^{2\lambda _s(\phi )}\mathrm{\Lambda }_s=0.`$ (A.2)
In (A.1)

$`T_{MN}=h_{\alpha \beta }\left(_M\phi ^\alpha _N\phi ^\beta {\displaystyle \frac{1}{2}}g_{MN}_P\phi ^\alpha ^P\phi ^\beta \right)V_\phi g_{MN}.`$ (A.3)
### B Power-law expansion

Here we derive the solution (4.3)-(4.6) from the general one. We start from the relation (3.3) that can be written in our case as follows

$$x^A(u)=\underset{sS}{}h_sU^{sA}\mathrm{ln}(\sqrt{|2\mathrm{\Lambda }_s/h_s|}u)+\overline{c}^A,$$
(B.4)
where $`u>0`$.

Introducing a new variable $`\tau >0`$ by formula

$$u=C\tau ^{1/(B+1)},C>0,$$
(B.5)
we rewrite (B.4) in the following manner

$$x^A=\underset{sS}{}h_sU^{sA}\mathrm{ln}\tau +x_0^A$$
(B.6)
with constants

$$x_0^A=\overline{c}^A\underset{sS}{}h_sU^{sA}\mathrm{ln}\left(\sqrt{|2\mathrm{\Lambda }_s/h_s|}C\right).$$
(B.7)
Due to orthogonality of $`U^s`$-vectors the constraints on integration constant $`U_A^s\overline{c}^A=0`$ may be written in the equivalent form

$$U_A^sx_0^A=\mathrm{ln}(\sqrt{|2\mathrm{\Lambda }_s/h_s|}C).$$
(B.8)
In components the solution (B.6) for $`(x^A)=(\varphi ^i,\phi ^\alpha )`$ reads

$`\varphi ^i={\displaystyle \frac{B}{(B+1)(D1)}}\mathrm{ln}\tau +\varphi _0^i,`$ (B.9)
$`\phi ^\alpha ={\displaystyle \frac{1}{B+1}}{\displaystyle \underset{sS}{}}h_s\lambda _s^\alpha \mathrm{ln}\tau +\phi _0^\alpha .`$ (B.10)
Here $`(x_0^A)=(\varphi _0^i,\phi _0^\alpha )`$.

For $`\gamma _0(\varphi )=\underset{i=0}{\overset{n}{}}d_i\varphi ^i`$ we get

$$\gamma _0(\varphi )=\frac{B}{(B+1)}\mathrm{ln}\tau +\varphi _0,$$
(B.11)
where $`\varphi _0=\underset{i=0}{\overset{n}{}}d_i\varphi _0^i`$.

From the definition of ”synchronous” variable

$$\mathrm{exp}(2\gamma _0(\varphi ))du^2=d\tau ^2$$
(B.12)
we obtain

$$C=|B+1|\mathrm{exp}(\varphi _0)$$
(B.13)
and hence (B.8) reads

$$\lambda _{s\alpha }\phi _0^\alpha =\mathrm{ln}(\sqrt{|2\mathrm{\Lambda }_s/h_s|}C).$$
(B.14)
The solution (4.3)-(4.6) follows from the formulas (B.9), (B.10), (B.12), (B.14) and $`A_i=\mathrm{exp}(2\varphi _0^i)`$.

### C Exponential expansion

The solution (4.10)-(4.12) may be obtained just along a line as it was done for the power-law case. The only difference here is the relation

$$u=C\mathrm{exp}(M\tau ),C>0$$
(C.15)
instead of (B.5). Using a procedure analogous to considered hereabove we get

$`\varphi ^i={\displaystyle \frac{1}{D1}}M\tau +\varphi _0^i,`$ (C.16)
$`\phi ^\alpha ={\displaystyle \underset{sS}{}}h_s\lambda _s^\alpha M\tau +\phi _0^\alpha .`$ (C.17)
and

$$\gamma _0(\varphi )=M\tau +\varphi _0.$$
(C.18)
From (B.12) we get

$$C=|M|^1\mathrm{exp}(\varphi _0)$$
(C.19)
and introducing new parameter $`m=M/(D1)`$ we obtain the solution (4.10)-(4.12).

### D $`b`$proper set of vectors

The statements of this subsection are based on relations (1.3) and (1.4)), that are used in Definition of $`b`$-proper set (see below). The conditions for exponential expansion follow from Theorem. Lemma 2 is useful for analysis of ”power-law” expansion.

In what follows $`\{\lambda _i\}`$ means a set of vectors $`\lambda _1,\mathrm{},\lambda _m\text{IR}^l,l\text{IN}`$. The vector space $`\text{IR}^l`$ is equipped with a scalar product (it corresponds to the positive definite matrix $`(h_{\alpha \beta })`$ in the bulk of the article).

Definition. Let $`b>0`$. The set of vectors $`\{\lambda _i\}`$ is called a $`b`$proper one if

$`\lambda _i\lambda _j=b,`$ (D.20)
$`\lambda _i^2b,`$ (D.21)
$`ij,i,j=1,\mathrm{}m.`$

Lemma 1. Let $`K=(\lambda _i\lambda _j)`$ be a matrix of scalar products for a $`b`$proper set of vectors $`\{\lambda _i\}`$. Then

$$\mathrm{det}K=(B+1)\underset{i=1}{\overset{m}{}}(\lambda _i^2b),$$
(D.22)
where

$$B=B(\lambda _1,\mathrm{},\lambda _m)\underset{i=1}{\overset{m}{}}\frac{b}{\lambda _i^2b}.$$
(D.23)
*Proof.* Using relations

$$K_{ij}=\lambda _i\lambda _j=(\lambda _i^2b)\delta _{ij}+b=(\lambda _i^2b)\left(\delta _{ij}+\frac{b}{\lambda _i^2b}\right)$$
(D.24)
we represent the matrix $`K`$ as a product of a matrix $`\overline{K}`$ and a diagonal matrix $`D`$:

$`\overline{K}_{ij}=\delta _{ij}+{\displaystyle \frac{b}{\lambda _i^2b}},D_{ij}=\delta _{ij}(\lambda _i^2b).`$ (D.25)
A product of determinants $`\mathrm{det}D`$ and $`\mathrm{det}\overline{K}`$ yields (D.22). Here we used the relation $`\mathrm{det}\overline{K}=B+1`$, that can be readily proved.

Theorem. Let $`\{\lambda _i\}`$ be a $`b`$proper set. The set $`\{\lambda _i\}`$ is linearly dependent if, and only if

$$B=B(\lambda _1,\mathrm{},\lambda _m)=1.$$
(D.26)
In this case

$$\underset{i=1}{\overset{m}{}}\frac{\lambda _i}{\lambda _i^2b}=0$$
(D.27)
and any subset of $`m1`$ vectors from $`\{\lambda _i\}`$ is linearly independent one.

*Proof.* The linear dependence of vectors $`\{\lambda _i\}`$ is equivalent to $`detK=0`$ that is equivalent (due to Lemma 1) to relation (D.26).

From the relation

$$B(\lambda _1,\mathrm{},\lambda _m)=B(\lambda _1,\mathrm{},\lambda _{m1})+\frac{b}{\lambda _m^2b}=1$$
(D.28)
we get $`B(\lambda _1,\mathrm{},\lambda _{m1})1`$ and hence the vectors $`\lambda _1,\mathrm{},\lambda _{m1}`$ are linearly independent. It is obvious that any $`m1`$ $`\lambda `$-vectors are also linearly independent.

To prove (D.27) let us denote

$$\lambda =\underset{i=1}{\overset{m}{}}\frac{\lambda _i}{\lambda _i^2b}.$$
(D.29)
Using the scalar products (D.24) we obtain

$$\lambda ^2=B(B+1)/b=0$$
(D.30)
and, hence, $`\lambda =0`$. The Theorem is proved.

As consequence we obtain the inequality on the number of $`b`$-proper vectors: $`ml+1`$. (The equality $`m=l+1`$ takes place when the vectors are linearly dependent.)

Lemma 2. In the $`b`$proper set $`\{\lambda _i\}`$:

$`a.`$ there are no two or more vectors $`\lambda _i`$ obeying $`\lambda _i^2<b`$;

$`b.`$ $`B>1`$ vectors $`\lambda _i`$ are linearly independent and all $`\lambda _i^2>b`$.

*Proof.* $`a.`$ Let us suppose that there are two vectors $`\lambda _i,\lambda _j`$, such that $`\lambda _i^2<b`$ and $`\lambda _j^2<b`$. Then $`b^2=(\lambda _i\lambda _j)^2\lambda _i^2\lambda _j^2<b^2`$, i.e. we come to a contradiction.

$`b.`$ If $`B>1`$ then by the Theorem the vectors are linearly independent and hence the matrix $`(\lambda _i\lambda _j)=(K_{ij})`$ is positive definite. This implies $`\mathrm{det}K>0`$. Due to relation (D.22) and part $`a.`$ of this lemma we get: $`\lambda _i^2>b`$ for all vectors.

Now, let the vectors $`\{\lambda _i\}`$ be linearly independent and all $`\lambda _i^2>b`$. Then $`detK>0`$, and due to (D.22), $`B+1>0`$. The lemma is proved.


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# 1 Introduction

## 1 Introduction

There has recently been much interest in a Hamilton-Jacobi (HJ) formulation for constrained systems -. Although there exists consensus in the formulation for purely first class systems, the case of second-class systems has to be dealt with separately.

For an unconstrained system described by a Lagrangian $`L(q_i,\dot{q}_i)`$, $`i=1,\mathrm{},n`$, the HJ-equations read,

$$H_0^{}p_0+H_0(q_0,q_i,p_i)=0$$
(1)
where $`H_0`$ is the canonical Hamiltonian, and where, for later convenience, $`p_0`$ and $`p_i`$ are shorthand for the partial derivatives

$$p_0=\frac{S}{q_0},p_i=\frac{S}{q_i},q_0=t,$$
(2)
with $`S`$ the Hamilton principal function $`S=S(t,q_i)`$ to be obtained as a solution of eq. (1).

Consider now a constrained system in which not all canonical variables are independent. In that case, the Lagrangian $`L(q_i,\dot{q}_i)`$ is singular, so that the determinant of the Hessian matrix $`H_{ij}=\frac{^2L}{\dot{q}_i\dot{q}_j}`$ is zero and the accelerations of some variables $`q_i`$ are not uniquely determined by the positions and the velocities at a given time.

Let $`nm_1`$ be the rank of the Hessian. In that case only $`nm_1`$ velocities $`\dot{q}_a`$ ($`a=m_1+1,\mathrm{},n`$) can be solved for as a function of the coordinates $`q_i`$, the momenta $`p_a`$ canonically conjugate to $`q_a`$, and the remaining velocities $`\dot{q}_\alpha `$, to yield $`\dot{q}_a=\dot{q}_a(q_i,p_b,\dot{q}_\beta )`$. The remaining momenta $`p_\alpha `$ ($`\alpha =1,\mathrm{},m_1`$) can be shown to be functions of $`q_i`$ and $`p_a`$ only, and represent the primary constraints $`\varphi _\alpha (q,p)p_\alpha +H_\alpha (q_i,p_a)=0`$, in the Dirac terminology . We write them in a form appropriate for our purposes:

$$H_\alpha ^{}p_\alpha +H_\alpha (q_0,q_i,p_a)=0,(\alpha =1,\mathrm{}m_1)$$
(3)
Adjoining these constraints to (1), we are led to consider the coupled set of differential equations

$`H_0^{}{\displaystyle \frac{S}{q_0}}+H_0(q_0,q_\alpha ,q_a,{\displaystyle \frac{S}{q_a}})=0,`$
$`H_\alpha ^{}{\displaystyle \frac{S}{q_\alpha }}+H_\alpha (q_0,q_\alpha ,q_a,{\displaystyle \frac{S}{q_a}})=0,\alpha =1,\mathrm{},m_1,`$ (4)
where $`H_0`$ is the canonical Hamiltonian evaluated on the subspace of the primary constraints, which only depends on the variables $`q_i`$, $`p_a`$ (and possibly $`q_0=t`$), as is well known. Continuing the above line of reasoning, the system of differential equations (1) should further be suplemented with the corresponding set of equations associated with possible secondary constraints $`\phi (q,p)_\sigma =0`$. Let there be $`m_2`$ such constraints. Adjoining these to (1), we have in our generalized notation,

$$H_{\underset{¯}{\alpha }}^{}(q_{\underset{¯}{i}},\frac{S}{q_{\underset{¯}{i}}})=0,\underset{¯}{i}=0,1,\mathrm{},n,\underset{¯}{\alpha }=0,1,\mathrm{},m=m_1+m_2.$$
(5)
Assuming that the secondary constraints do not represent constraints among the coordinates $`q_{\underset{¯}{i}}`$ themselves, 333As we shall see, this condition is not satisfied for the “non-linear sigma model”. this coupled set of differential equations only admits a solution provided the $`H_{\underset{¯}{\alpha }}^{}`$ are in weak involution in the sense 444On the space of solutions to (1) this then implies strong involution .

$$\{H_{\underset{¯}{\alpha }}^{},H_{\underset{¯}{\beta }}^{}\}|_{p_{\underset{¯}{i}}=\frac{S}{q_{\underset{¯}{i}}}}=C_{\underset{¯}{\alpha },\underset{¯}{\beta }}^{\underset{¯}{\gamma }}H_{\underset{¯}{\gamma }}^{}(q_{\underset{¯}{i}},\frac{S}{q_{\underset{¯}{i}}}),$$
(6)
where the (generalized) Poisson bracket is defined by

$$\{A,B\}=\underset{\underset{¯}{i}=0}{\overset{m}{}}\left(\frac{A}{q_{\underset{¯}{i}}}\frac{B}{p_{\underset{¯}{i}}}\frac{B}{q_{\underset{¯}{i}}}\frac{A}{p_{\underset{¯}{i}}}\right).$$
(7)
Hence, for second class systems, our Ansatz (5) for the HJ equations must be incorrect.

One way of circumventing the integrability problem would be to embed the model in question into a larger phase space, such as to become first class . Then the integrability conditions (6) are statisfied and our Ansatz is self-consistent . This is, however, not the point of view we wish to take. As we shall rather illustrate in terms of a number of examples, a direct way of proceeding in the case of systems with second class constraints is to perform a canonical change of variables in which the second-class constraints become part of the variables. If the Poisson brackets of constraints are given by constant matrices, one can always perform a transformation to variables, where the constraints themselves can be grouped into canonically conjugate pairs. Otherwise this may not be possible , and the embedding procedure would provide a way out of the dilemma.

Unlike the discussions found in the literature  it is the final objective of this paper to actually construct the Hamilton-principal function $`S`$ for the examples we consider, and to obtain from there the solution of the corresponding Euler-Lagrange equations following standard methods.

Our strategy to be illustrated by a number of examples will be the following: We begin by seeking all constraints, following the Dirac algorithm. We then seek a canonical transformation from a set of phase space variables $`q_i,p_i`$ to a new set $`(Q,P)`$ = ($`q_a^{},\chi _\alpha ,p_a^{},𝒫_\alpha `$), in which the complete set of primary and secondary constraints are grouped into canonically conjugate pairs ($`\chi _\alpha ,𝒫_\alpha `$) satisfying $`\{\chi _\alpha ,𝒫_\beta \}=\delta _{\alpha \beta }`$. Since the number of second class constraints is even, this is always possible for linear constraints. 555As shown in ref.  this is not always possible, if the constraints are nonlinear. The remaining set of variables $`(q_a^{},p_a^{})`$, which in turn are canonically conjugate to each other, $`\{q_a^{},p_b^{}\}=\delta _{ab}`$, are chosen to commute with the set ($`\chi _\alpha ,𝒫_\alpha `$). Denoting by $`F_2(q;p^{},𝒫;t)`$ the generating function for this canonical transformation, the new Hamiltonian will be given by $`\stackrel{~}{H}(q^{},p^{};\chi ,𝒫;t)=H(q(q^{},p^{},\chi ,𝒫),p(q^{},p^{},\chi ,𝒫);t)+\frac{F_2}{t}`$, where $`q(q^{},p^{},\chi ,𝒫)`$ and $`p(q^{},p^{},\chi ,𝒫)`$ are obtained by solving the coupled set of equations $`p_i=\frac{F_2}{q_i}`$, $`\chi _\alpha =\frac{F_2}{𝒫_\alpha }`$ and $`q_a^{}=\frac{F_2}{p_a^{}}`$. The extended action in terms of the new variables now reads, 666To formulate the general problem it is convenient to work with the extended action, where all of the constraints are implemented via Lagrange multipliers. Since all the constraints are supposed to be generated from the total Hamiltonian $`H_T`$, involving only the primary constraints, via the Dirac self-consistent algorithm, the Lagrange multipliers associated with the secondary constraints will eventually turn out to be zero in a variational calculation.

$$S_E=𝑑t\{p_a^{}\dot{q}_a^{}+𝒫_\alpha \dot{\chi }_\alpha \stackrel{~}{H}(q^{},p^{},\chi ,𝒫;t)\lambda _\alpha \chi _\alpha \eta _\alpha 𝒫_\alpha \}$$
(8)
Since the constraints commute with $`q_a^{}`$ and $`p_a^{}`$, the equations of motion for these variables are determined by $`\stackrel{~}{H}`$ with the constraints set equal to zero, strongly:

$$\dot{q}_a^{}=\{q_a^{},\stackrel{~}{H}(q^{},p^{},0,0;t)\},\dot{p}_a^{}=\{p_a^{},\stackrel{~}{H}(q^{},p^{},0,0;t)\}$$

while the equations for the remaining coordinates, the constraints, just reproduce the persistency equations, and represent a set of differential equations in the new coordinates. With respect to the coordinates $`(q^{},p^{})`$ the problem has thus been reduced to that of an unconstrained system, and the usual Hamilton-Jacobi treatment applies. Note that the canonical change in phase space variables has now removed the inconsisteny of the originally envisaged HJ formulation. The reason is that, unlike the Hamiltonian, the Hamilton Principal function in the new coordinates $`(Q,P)`$ cannot be obtained by a simple substitution $`(q,p)(q(Q,P),p(Q,P))`$.

In different terms: from the total action (8) we obtain for the HJ equation

$$\stackrel{~}{H}(q^{},\frac{S}{q^{}},\chi ,\frac{S}{\chi },t)+\frac{S}{t}=0$$
(9)
together with the constraint equations,

$$\chi _\alpha =0,𝒫_\alpha =\frac{S}{\chi _\alpha }=0.$$
(10)
In general our problem thus reduces to solving the differential equations

$$\widehat{H}(q^{},\frac{S}{q^{}},t)+\frac{S}{t}=0$$
(11)
with

$$\widehat{H}(q^{},p^{})=\stackrel{~}{H}(q^{},p^{},\chi =𝒫=0).$$
(12)
In the following section we illustrate these ideas in terms of a number of simple examples. We shall be rather detailed in the first example, in order to illustrate various aspects of the problem.

## 2 Second class systems and HJ formulation

In this section we shall consider several simple quantum mechanical models of purely second class systems and show how to properly deal with them in order to construct the corresponding Hamilton Principal function $`S`$. We also verify that the formalism yields the solutions to the corresponding equations of motion.

Example 1

Consider the Lagrangean

$$L=\frac{1}{2}\dot{x}^2+\dot{x}y\frac{1}{2}(xy)^2$$
(13)
The corresponding Euler-Lagrange equations can be written in the form, $`\ddot{x}=0,\dot{x}+x=y`$. We have one primary constraint $`\varphi :=p_y=0`$, and the canonical Hamiltonian on the primary surface reads,

$$H_0=\frac{1}{2}(p_xy)^2+\frac{1}{2}(xy)^2.$$
(14)
The usual Dirac algorithm yields a secondary constraint $`\phi =0`$,

$$\{\varphi ,H_0\}=2\phi ,\phi =y\frac{1}{2}(p_x+x).$$
(15)
There are no further constraints. Since $`\{\phi ,\varphi \}=1`$ these constraints are second class and canonically conjugate to each other.

If we were to proceed naivly and make the replacements $`p_y=\frac{S}{y}`$ in the primary constraint $`p_y=0`$, then this would imply that $`S`$ is independent of $`y`$, which would be in direct conflict with the equation obtained by making the substitution $`p_x=\frac{S}{x}`$ in the constraint $`\phi =0`$. This reflects the non-commutative structure of our second class system, as well as the violation of the integrability condition (6).

In order to circumvent this difficulty, we make a canonical transformation to a new set of variables $`(q^{},p^{})`$ and $`(\chi ,𝒫)`$, in which one of the canonically conjugate pairs are chosen to be the constraints themselves:

$`\chi `$ $`=`$ $`\phi ,𝒫=\varphi ,`$
$`q^{}`$ $`=`$ $`x{\displaystyle \frac{1}{2}}p_y,p^{}=p_x+{\displaystyle \frac{1}{2}}p_y.`$ (16)
We evidently have

$$\{\chi ,𝒫\}=1,\{q^{},p^{}\}=1.$$
(17)
Futhermore, $`(q^{},p^{})`$ and $`(\chi ,𝒫)`$ “commute”. Our change of variables is thus canonical. Indeed one readily checks that

$$p_x\dot{x}+p_y\dot{y}=p^{}\dot{q}^{}+𝒫\dot{\chi }+\frac{1}{2}\frac{d}{dt}(p^{}𝒫\frac{1}{2}𝒫^2).$$

Furthermore, the corresponding generating functional $`F_2(x,y;p^{},𝒫)`$ for the transformation (2) is readily constructed:

$$F_2=(p^{}\frac{1}{2}𝒫)x+𝒫y\frac{1}{2}𝒫p^{}+\frac{1}{8}𝒫^2.$$

The inverse transformation to (2) reads

$`x`$ $`=`$ $`q^{}+{\displaystyle \frac{1}{2}}𝒫,p_x=p^{}{\displaystyle \frac{1}{2}}𝒫,`$
$`y`$ $`=`$ $`{\displaystyle \frac{1}{2}}(q^{}+p^{})+\chi ,p_y=𝒫.`$ (18)
In terms of the new variables,

$$\stackrel{~}{H}=\frac{1}{2}\left[\frac{1}{2}(p^{}q^{})(\frac{1}{2}𝒫+\chi )\right]^2+\frac{1}{2}\left[\frac{1}{2}(q^{}p^{})(\frac{1}{2}𝒫\chi )\right]^2$$
(19)
One checks,

$`\dot{q}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p^{}q^{}){\displaystyle \frac{1}{2}}𝒫`$
$`\dot{p}^{}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(p^{}q^{}){\displaystyle \frac{1}{2}}𝒫`$
$`\dot{𝒫}`$ $`=`$ $`2\chi `$
$`\dot{\chi }`$ $`=`$ $`\eta {\displaystyle \frac{1}{2}}(p^{}q^{})+{\displaystyle \frac{1}{2}}𝒫`$
The 3’rd equation reproduces the secondary constraint, and the 4’th equation fixes the parameter $`\eta `$ as expected.

The above shows, that as a consequence of having the conjugate pairs in strong involution, we can simply set the constraints in $`\stackrel{~}{H}_0`$ equal to zero and proceed with the Hamiltonian

$$\widehat{H}_0=\frac{1}{4}(p^{}q^{})^2.$$
(21)
Hence we are left with only one HJ equation,

$$\frac{S}{t}+\frac{1}{4}(\frac{S}{q^{}}q^{})^2=0$$
(22)
showing that $`S`$ will be just a function of $`q^{}`$. The HJ equation associated with the primary constraint $`𝒫=0`$ tells us that $`S`$ is independent of the new coordinate $`\chi `$: $`𝒫=\frac{S}{\chi }=0`$.

Since the Hamiltonian does not depend explicitely on time, we must have

$$\frac{S}{t}=\alpha ^2=constant$$

so that equation (22) has the solution

$$S(q^{},\alpha ,t)=\pm 2\alpha q^{}+\frac{1}{2}q^2\alpha ^2t.$$
(23)
Notice that $`S`$ depends only on one integration constant, since when regarded as the generating function taking us to new momenta $`p^{}`$ and $`𝒫`$, one of the momenta ($`𝒫`$) is just the primary constraint $`𝒫=\varphi =0`$. There is another constant arising from the equation

$$\beta =\frac{S}{\alpha },$$

which implies from (23),

$$\beta =\pm 2q^{}2\alpha t.$$

Together with (we choose the plus sign in (23), without loss of generality)

$$p^{}=\frac{S}{q^{}}=2\alpha +q^{}$$

we now solve the problem in the standard way. Taking account of the constraints $`\chi _\alpha =𝒫_\alpha =0`$ in (2), we find

$$x=\frac{1}{2}\beta +\alpha t,y=(\alpha +\frac{1}{2}\beta )+\alpha t$$

which evidently is a solution of the Euler-Lagrange equations associated with (13).

Example 2: QM $`\sigma `$ \- model

Consider now the quantum mechanical version of the “sigma model” as given by the Lagrangean

$$L=\frac{\dot{\stackrel{}{q}}^2}{2}+\lambda (\stackrel{}{q}^21)$$
(24)
implying the primary constraint $`p_\lambda =0`$. The canonical Hamiltonian on the primary surface reads,

$$H_0=\frac{\stackrel{}{p}^2}{2}\lambda (\stackrel{}{q}^21).$$
(25)
The complete set of primary and secondary constraints can be written in the form

$$\varphi =p_\lambda ,\phi =\lambda +\frac{\stackrel{}{p}^2}{2},\phi _1=\stackrel{}{q}^21,\phi _2=\frac{\stackrel{}{p}\stackrel{}{q}}{2\stackrel{}{q}^2}.$$
(26)
The Poisson brackets of the constraints show that they can be grouped into mutually commuting, canonically conjugate pairs:

$$\{\phi ,\varphi \}=1,\{\phi _1,\phi _2\}=1.$$
(27)
Let us consider for simplicity the case of two dimensions. One may be tempted to make a transformation from $`(\lambda ,p_\lambda )`$ to $`(\phi ,\varphi )`$ and $`(q_1,q_2,p_1,p_2)`$ to $`(\phi _1,\phi _2,\theta ,J)`$, where ($`\theta ,J`$) are the canonically conjugate variables,

$$\theta =\mathrm{arctan}(q_2/q_1),J=ϵ_{ij}q_ip_j,$$
(28)
with $`J`$ evidently playing the role of the “third” component of angular momentum. The second set of variables commute with the first set, except for $`\theta `$, which does not commute with $`\phi `$: $`\{\phi ,\theta \}=J/\stackrel{}{q}^2`$. Hence this transformation is not canonical.

In order to remove this problem we eliminate $`\lambda `$ and $`p_\lambda `$ using the constraints $`\phi =0`$ and $`\varphi =p_\lambda =0`$. 777It is easy to see, that this is a legitimate procedure since $`\lambda `$ plays the role of a Lagrange multiplier in the Lagrangean. More precisely, implementing strongly the constraint $`\phi =0`$ in the Hamiltonean (25) from the outset, leads to an equivalent set of equations of motion. The resulting Hamiltonian is $`H_0=\frac{1}{2}\stackrel{}{p}^2q^2`$. Noting that

$$J^2=\stackrel{}{q}^2\stackrel{}{p}^2(\stackrel{}{q}\stackrel{}{p})^2=\stackrel{}{p}^2(1+\phi _1)4\phi _2^2(\phi _1+1)^2$$

we then obtain for $`\widehat{H}_0`$,

$$\widehat{H}_0=\frac{1}{2}J^2$$

We are thus left to solve the equation

$$\frac{S}{t}+\frac{1}{2}\left(\frac{S}{\theta }\right)^2=0$$
(29)
which has the solution (an irrelevant additive constant has been ommitted)

$$S(\theta ,\alpha ,t)=\frac{\alpha ^2}{2}t+\alpha \theta .$$
(30)
In order to obtain the solution to the Euler-Lagrange equations we recall that we still have the equation $`\frac{S}{\alpha }=\beta `$, which gives the solution in terms of two integrations constants:

$$\theta =\alpha t+\beta $$
(31)
One again verifies that this provides the solution of the Euler-Lagrange equations.

Example 3: Multidimensional rotator

Consider the Lagrangean ,

$$L=\frac{1}{2}\dot{\stackrel{}{q}}^2+\lambda \stackrel{}{q}\dot{\stackrel{}{q}}.$$

There is one primary constraint $`\varphi =p_\lambda =0`$ and the canonical Hamiltonian evaluated on the primary surface reads,

$$H_0=\frac{1}{2}(\stackrel{}{p}\lambda \stackrel{}{q})^2.$$
(32)
The model effectively describes the motion on a $`n1`$ dimensional sphere without specification of the radius of the sphere. Note that unlike in the case of the $`\sigma `$-model, the momentum conjugate to $`\stackrel{}{q}`$ is no longer the mecanical momentum: $`\stackrel{}{p}=\dot{\stackrel{}{q}}+\lambda \stackrel{}{q}`$. There is only one secondary constraint

$$\phi =\lambda \frac{\stackrel{}{p}\stackrel{}{q}}{\stackrel{}{q}^2},$$
(33)
which has been chosen to be canonically conjugate to $`\varphi `$: $`\{\phi ,\varphi \}=1`$.

i) We first consider this model in two dimensions. A coordinate change $`(\lambda ,p_\lambda )(\phi ,\varphi )`$ and $`(q_1,p_1,q_2,p_2)(r,p_r,\theta ,J)`$ would again lead to the second set not commuting with the first, so that $`\lambda `$ should first be eliminated using the constraint $`\phi =0`$. Having done this we may replace the canonical Hamiltonian (32) by (see footnote 7)

$$\widehat{H}_0=\frac{1}{2}\left(\stackrel{}{p}\frac{(\stackrel{}{p}\stackrel{}{q})}{\stackrel{}{q}^2}\stackrel{}{q}\right)^2=\frac{J^2}{2r^2}$$
(34)
where $`r=\sqrt{\stackrel{}{q}^2}`$, $`J`$ is given in (28), and $`\{\theta ,J\}=1`$. Notice that $`r`$ just enters as a parameter. Replacing $`J`$ by $`\left(\frac{S}{\theta }\right)`$, we obtain from (34) the HJ equation,

$$\frac{S}{t}+\frac{1}{2r^2}\left(\frac{S}{\theta }\right)^2=0,$$
(35)
which just differs from (29) by $`r`$ being arbitrary, and not equal to one. Correspondingly (30) is replaced by

$$S(\theta ,\alpha ,t)=\frac{\alpha ^2}{2}t+r\alpha \theta .$$
(36)
The rest proceeds as in the case of model 2.

ii) We now reconsider the problem in $`n`$-dimensions. Solving again the constraint $`\phi =0`$ for $`\lambda `$, the HJ equation (35) is replaced by, 888Note that the coordinates $`q_a`$ play the role of $`q_a^{}`$ in section 1.

$$\frac{S}{t}+\frac{1}{2}\left(\frac{S}{q_a}\right)^2\frac{1}{2}\left(\frac{S}{q_a}\right)\left(\frac{S}{q_b}\right)\frac{q_aq_b}{\stackrel{}{q}^2}=0.$$

We make the Ansatz,

$$S=f(\stackrel{}{n}\stackrel{}{q})$$
(37)
with $`\stackrel{}{n}`$ a unit normal vector parametrized by $`n1`$ constants. The HJ equation for $`f`$ then reads,

$$\frac{1}{2}\left(1\frac{(\stackrel{}{n}\stackrel{}{q})^2}{\stackrel{}{q}^2}\right)f^2(\stackrel{}{n}\stackrel{}{q})=\frac{\alpha ^2}{2}.$$

Setting $`x=\stackrel{}{n}\stackrel{}{q}`$ and $`r^2=\stackrel{}{q}^2`$, we then find

$$f^{}(x)=\pm \frac{\alpha }{\sqrt{1\frac{x^2}{r^2}}}$$

so that upon integration in $`x`$ the Hamilton principal function takes the form

$$S=\alpha r\mathrm{tan}^1\frac{\stackrel{}{n}\stackrel{}{q}}{\sqrt{r^2(\stackrel{}{n}\stackrel{}{q})^2}}\frac{\alpha ^2}{2}t+const$$
(38)
The Hamilton principal function contains $`n`$ independent constants, which we take to be $`\alpha `$ and $`n_1,n_2,\mathrm{}n_{n1}`$, while the normalization of $`\stackrel{}{n}`$ implies for the $`n`$’th component, $`n_n=\sqrt{1_{a=1}^{n1}n_a^2}`$. Differentiating the Hamilton principal function with respect to these constants (new momenta in the corresponding generating functional) yields the $`n`$ time-independent new coordinates:

$$\beta =\frac{S}{\alpha },\beta _a=\frac{S}{n_a},a=1,2,\mathrm{}n1.$$
(39)
From the first equation and (38) we easily obtain

$$\stackrel{}{n}\stackrel{}{q}=r\mathrm{cos}\frac{\beta +\alpha t}{r}r\mathrm{cos}\mathrm{\Omega }(t).$$
(40)
Using this, we obtain from the second equation in (39)

$$q_a=\frac{\beta _a}{\alpha }\mathrm{sin}\mathrm{\Omega }(t)+\frac{n_aq_n}{n_n};a=1,2,\mathrm{},n1$$
(41)
Multiplying this equation by $`n_a`$, summing from $`a=1`$ to $`a=n1`$ and using the normalization of $`\stackrel{}{n}`$ we obtain from (41),

$$\stackrel{}{n}\stackrel{}{q}=\frac{1}{\alpha }\underset{a=1}{\overset{n1}{}}n_a\beta _a\mathrm{sin}\mathrm{\Omega }(t)+\frac{q_n}{n_n}=r\mathrm{cos}\mathrm{\Omega }(t)$$

This allows us to solve for $`q_n`$ as

$$q_n=n_n\left(r\mathrm{cos}\mathrm{\Omega }(t)\frac{1}{\alpha }\underset{a=1}{\overset{n1}{}}n_a\beta _a\mathrm{sin}\mathrm{\Omega }(t)\right).$$

Substituting back in the equation for $`q_a`$ we obtain

$$q_a=\frac{1}{\alpha }\left(\beta _an_a\underset{a=1}{\overset{n1}{}}n_a\beta _a\right)\mathrm{sin}\mathrm{\Omega }(t)+n_ar\mathrm{cos}\mathrm{\Omega }(t),a=1,2,\mathrm{},n1.$$

Introducing the $`n`$-dimensional vector $`\stackrel{}{\beta }=(\beta _1,\beta _2,\mathrm{},\beta _{n1},0)`$, these results can be written compactly as

$$\stackrel{}{q}=\frac{1}{\alpha }\left(\stackrel{}{\beta }(\stackrel{}{n}\stackrel{}{\beta })\stackrel{}{n}\right)\mathrm{sin}\mathrm{\Omega }(t)+r\stackrel{}{n}\mathrm{cos}\mathrm{\Omega }(t)$$

Substituting the above result into our original condition, $`r^2=\stackrel{}{q}^2`$, leads to

$$r=\sqrt{\frac{\stackrel{}{\beta }^2(\stackrel{}{n}\stackrel{}{\beta })^2}{\alpha ^2}}$$
(42)
Thus the radius of motion is fixed if the new time independent coordinates are specified. More conventionally one would consider the radius of motion as a given initial condition; then the above condition eliminates one of the $`\beta _a`$ as an independent constant. Indeed, from (42) we see that the component of $`\stackrel{}{\beta }`$ orthogonal to $`\stackrel{}{n}`$ is fixed to be $`\alpha r`$. In this case one therefore considers $`\alpha `$, $`r`$ and the remaining $`n2`$ constants, which determine the component of $`\stackrel{}{\beta }`$ parallel to $`\stackrel{}{n}`$ as independent constants.

Example 4: Landau model in the zero mass limit

The Lagrangean of a spinless charged particle moving on a two dimensional plane in a constant background magnetic field $`B`$ perpendicular to the 12-plane, and a harmonic oscillator potential, can be written in the form

$$L_{\mathrm{𝐿𝑎𝑛𝑑𝑎𝑢}}=\frac{m}{2}\dot{\stackrel{}{q}}^2+\frac{B}{2}\stackrel{}{q}\times \dot{\stackrel{}{q}}\frac{k}{2}\stackrel{}{q}^2.$$

In the zero-mass limit the Lagrangian of the model reduces to

$$L=\frac{B}{2}\stackrel{}{q}\times \dot{\stackrel{}{q}}\frac{k}{2}\stackrel{}{q}^2$$
(43)
There are only two primary constraints

$$\varphi _i=\frac{1}{\sqrt{B}}p_i+\frac{\sqrt{B}}{2}ϵ_{ij}q_j$$

with non-vanishing Poisson brackets

$$\{\varphi _i,\varphi _j\}=ϵ_{ij}$$

We perform the following change of variables from $`q_i,p_i`$ to $`\varphi _1,\varphi _2,q,p`$: 999In the notation of section 1, $`q,p`$ now play the role of $`q^{},p^{}`$, and $`\varphi _1=\chi ,\varphi _2=𝒫`$.

$`\varphi _1`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{B}}}p_1+{\displaystyle \frac{\sqrt{B}}{2}}q_2,\varphi _2={\displaystyle \frac{1}{\sqrt{B}}}p_2{\displaystyle \frac{\sqrt{B}}{2}}q_1,`$
$`q`$ $`=`$ $`{\displaystyle \frac{1}{2}}q_1+{\displaystyle \frac{1}{B}}p_2,p=p_1{\displaystyle \frac{B}{2}}q_2,`$ (44)
with the inverse transformation

$`q_1`$ $`=`$ $`q{\displaystyle \frac{1}{\sqrt{B}}}\varphi _2,q_2={\displaystyle \frac{1}{B}}(\sqrt{B}\varphi _1p),`$
$`p_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\sqrt{B}\varphi _1+p),p_2={\displaystyle \frac{1}{2}}(Bq+\sqrt{B}\varphi _2).`$ (45)
In terms of the new coordinates the canonical Hamiltonian evaluated on the constraint surface reads,

$$\widehat{H}_0=\frac{k}{2}\left(q^2+\frac{1}{B^2}p^2\right)$$
(46)
Hence we have for the HJ equation

$$\frac{S}{t}+\frac{k}{2B^2}\left(B^2q^2+\left(\frac{S}{q}\right)\right)=0,$$
(47)
having the solution

$$S(q,\alpha ,t)=\alpha t+B^q𝑑q^{}\sqrt{\frac{2\alpha }{k}q^2}+const.$$
(48)
From the equation $`\beta =\frac{S}{\alpha }`$ we then obtain in the usual way the solution

$$q(t)=\sqrt{\frac{2\alpha }{k}}\mathrm{cos}(\omega t+b),$$

with $`\omega =k/B`$ and b = $`\beta k/B`$.

## 3 Conclusion

A naive extension of Hamilton-Jacobi (HJ) theory to second class constrained systems is condemned to failure right from the outset since it leads to differential equations which are in direct conflict. Mathematically this conflict is expressed by the violation of the integrability condition (6).

It is useful to reformulate these integrability conditions in terms of the commutator algebra of operators. To this end we begin by introducing for each $`H_{\underset{¯}{\alpha }}^{}`$ a linear differential operator $`X_{\underset{¯}{\alpha }}`$ defined by

$$X_{\underset{¯}{\alpha }}f=\{H_{\underset{¯}{\alpha }}^{},f\},$$
(49)
where the Poisson bracket is the generalized Poisson bracket defined in (7), and $`f`$ an arbitrary functions of $`q_{\underset{¯}{i}}`$ and $`p_{\underset{¯}{i}}`$. If $`H_{\underset{¯}{\alpha }}^{}`$ does not possess an explicit time dependence, we recover the usual definition of the Poisson bracket. Making use of the Jacobi identity for Poisson brackets, we obtain from Eq.(49) the following relation between the commutator of the operators $`X_{\underset{¯}{\alpha }}`$, and the Poisson algebra of the constraints :

$$[X_{\underset{¯}{\alpha }},X_{\underset{¯}{\beta }}]f=\{\{H_{\underset{¯}{\alpha }}^{},H_{\underset{¯}{\beta }}^{}\},f\}.$$
(50)
for arbitrary $`f`$. It is now clear that the commutator algebra of the operators $`X_{\underset{¯}{\alpha }}`$ will not close unless the Poisson algebra of the $`H_{\underset{¯}{\alpha }}^{}`$ closes.

One way in which this problem has been approached is to iteratively extend the set of operators $`H_\alpha ^{}`$ until the algebra closes. Equivalently, new operators in (50) are iteratively introduced, enhancing the original set $`X_{\underset{¯}{\alpha }}`$ to a larger set $`X_{\overline{\alpha }}`$ of operators until the closure relation

$$[X_{\overline{\alpha }},X_{\overline{\beta }}]=c_{\overline{\alpha }\overline{\beta }}^{\overline{\gamma }}X_{\overline{\gamma }}.$$
(51)
is achieved . The corresponding system of partial differential equations will then be integrable. However, there is no general reason why the solution for the Hamilton principal function $`S`$ for this extended first-class algebra should provide the solution for the problem in question.

As an alternative approach to formulate a consistent set of HJ equations, one can consider  the embedding of the second class system into a first class one characterized by a strongly involutive algebra, following the BFT  construction. This was not the point of view to be taken in this paper. It was rather our objective to confront ourselves directly with the problem of non-integrability, and to point out a procedure for coping with this problem. The procedure consisted in making a canonical transformation to a new set of variables in which the second class constraints become part of the variables, grouped in canonically conjugate pairs. As we illustrated by a number of examples, the new Hamiltonian is then obtained by setting the constraints strongly equal to zero, and the HJ equations associated with the constraints are no longer in conflict with each other, nor with this new Hamiltonian. The solution for the coordinates obtained in this way were shown to coincide with the solutions of the corresponding Lagrange equations of motion.

If the constraints are nonlinear, one is in general not able to perform a canonical transformation, turning these constraints into canonical pairs , and our analysis as presented here is no longer valid. Investigations in this directions are presently under way.

Acknowledgement

One of the authors (K.D.R) would like to thank the Department of Physics of the University of Stellenbosch for the warm hospitality extended to him during his stay, as well as Heinz J. Rothe for a useful discussion.


text:
# Generation of photon number states

## Abstract

The Visible Light Photon Counter (VLPC) has the capability to discriminate photon number states, in contrast to conventional photon counters which can only detect the presence or absence of photons. We use this capability, along with the process of parametric down-conversion, to generate photon number states. We experimentally demonstrate generation of states containing 1,2,3 and 4 photons with high fidelity. We then explore the effect the detection efficiency of the VLPC has on the generation rate and fidelity of the created states.

preprint: APS/123-QED

Photon number states play an important role in quantum optics. Such states exhibit effects which contradict classical electromagnetic theory, such as anti-bunching and negativity of the Wigner function Walls and Milburn (1995). They also have important practical applications. In optical telecommunication, they achieve an optimal channel capacity Yamamoto and Haus (1986); Caves and Drummond (1994). They can also improve the sensitivity of an interferometer to the Heisenberg limit Holland and Burnett (2001). In this limit, the minimum detectable phase shift is inversely proportional to the number of photons, rather then the square root of the number achieved by the standard quantum limit.

Recently, there has been tremendous effort in generating single photon states. This is often achieved by using single emitters such as molecules, quantum dots, or diamond color centers Santori et al. (2001); Kim et al. (1999a); Lounis and Moerner (2000); Michler et al. (2000); Moreau et al. (2001); Beveratos et al. (2002); Yuan et al. (2002). Although these sources hold great promise, they currently suffer from substantial collection losses. This creates a large vacuum contribution, which means that the generated state greatly deviates from one with exactly one photon. Furthermore, the extension of these devices to generating higher order photon number states remains a difficult problem. The generation of a two photon number state has been recently demonstrated in the microwave regime Varcoe et al. (2001), using Rydberg atoms in a high-Q cavity. There have also been proposals for generation higher order optical photon number states using trapped atoms in an optical cavity Brown et al. (2003). The implementation of these schemes remains experimentally challenging.

There are generally two ways to engineer the quantum state of a system. The first is to start the system in a well known initial state, and turn on a unitary time evolution which transforms the initial state to the desired state. The second method is to make a non-destructive measurement of the quantum system, and project it onto the desired Hilbert space using the projection postulate of quantum mechanics. In this letter, we use the latter method of projection and conditional post-selection to generate a photon number state. This is done using the non-linear process of parametric down-conversion (PDC) in conjunction with a photon number detector known as the Visible Light Photon Counter (VLPC). In PDC, a bright pump field is injected into a material which exhibits a second order non-linear dipole response. This non-linearity can cause a pump photon to spontaneously split into two photons of lower energy, referred to as the signal and idler. By appropriately adjusting the phase matching conditions of the process, the signal and idler can be made to propagate in different directions, allowing them to be spatially resolved. Since these two photons always come in pairs, if a signal photon is detected, there must be an idler photon in the conjugate mode. If a photon counter is placed in front of the signal arm and detects one photon, the corresponding idler arm is prepared in a state containing only one photon by the projection postulate of quantum mechanics. The photon counter in essence performs a non-destructive measurement of the idler arm.

The above scheme can be extended to generation of higher order number states. Suppose a short pump laser pulse is injected into a non-linear crystal. We can define the photon number operator $`\widehat{n}=_T\widehat{𝐚}^{}\left(t\right)\widehat{𝐚}\left(t\right)𝑑t`$, where the integral is taken over the pulse duration of the pump and $`\widehat{𝐚}\left(t\right)`$ is the bosonic photon annihilation operator in the time domain. If n signal photons are detected by a photon counter in a given pulse, then the idler pulse is projected onto a state $`|\psi _i`$ which satisfies the condition

$$\widehat{n}|\psi _i=n|\psi _i.$$
(1)
That is, the idler mode is an eigenstate of the number operator with eigenvalue $`n`$. States which satisfy this special property are known as photon number states.

In order to implement the above scheme, the photon counter must be able to determine the exact number of photons in the signal arm within the short time duration of the pump pulse. Conventional photon counters such as avalanche photo-diodes (APDs) cannot do this, because they suffer from long dead time and large multiplication noise. Such detectors can only distinguish the zero photon case from the non-zero photon case in a pulse, making it impossible to generate photon number states.

Recently, a new type of photon counter known as the Visible Light Photon Counter (VLPC) Turner et al. (1993), has been shown to have the capability to distinguish different photon number states with high quantum efficiency Takeuchi et al. (1999); Atac et al. (1994); Kim et al. (1999b); Waks et al. . When two photons are simultaneously detected by the VLPC, it generates an electrical pulse which is twice as large. This pattern holds for higher photon numbers. The pulse height can then be used to determine the number of detected photons.

The VLPC can do photon number detection due to two unique properties. First, the VLPC has an active area of 1mm in diameter. When one photon is detected, it forms a 5$`\mu `$m diameter dead spot on the detector surface, leaving the rest of the active area available for subsequent detection events. As long as the probability that two photons land on the same spot is small, which is true if the incident light in not too tightly focussed, all of the photons will be detected. This circumvents the dead time problem which conventional APDs suffer from. Second, the VLPC has extremely low multiplication noise. Multiplication noise refers to the fluctuations in the number of electrons emitted by a detector in response to a photodetection event. These fluctuations limit the ability of the detector to infer photon number Waks et al. . Multiplication noise is quantified by a parameter known as the Excess Noise Factor $`F=M^2/M^2`$, where $`M`$ is the number of electrons emitted by the detector due to a photon detection. In the ideal limit where $`F=1`$, the detector emits a deterministic number of electrons for each detection. APDs with single carrier multiplication, which achieves the best noise performance, still have a large excess noise factor of $`F=2`$ McIntyre (1966), making it impossible to discriminate photon numbers. The VLPC, in contrast, achieves nearly noise free multiplication with $`F=1.03`$ Kim et al. (1997); Waks et al. .

The experimental setup for generating photon number states is shown in Figure 1. A 266nm pump source is generated from the fourth harmonic of a Q-switched Nd:YAG laser. The pump pulses have a duration of 20ns, and a repetition rate of 45kHz. A dispersion prism separates the fourth harmonic from the residual second harmonic, which is used to illuminate a high speed photodiode to generate a triggering signal. The fourth harmonic pumps a BBO crystal, which is set for non-collinear degenerate phase-matching. In this condition, the signal and idler waves are both 532nm in wavelength and have a divergence angle of 1 degree from the pump. The pump is loosely focused before the BBO crystal to achieve a minimum waste at the collection lens. This results in a sharper two-photon image which enhances the collection efficiency Monken et al. (1998).

Two VLPC detectors are used in this experiment. Each detector is held in a separate helium bath cryostat and cooled down to 6-7K, which is the optimum operating temperature. The VLPC is sensitive to photons with wavelengths of up to 30$`\mu `$m, so it must be shielded from room temperature thermal radiation. This is achieved by encasing the detector in a copper shield, which is cooled down to 6K. Acrylic windows at the front of the copper shield are used as infra red filters. These windows are highly transparent at visible wavelengths and simultaneously nearly opaque at 2-30$`\mu `$m wavelengths.

VLPC 1 is used as the triggering detector which detects the number of photons generated in the signal arm on a given laser pulse. The output of VLPC 1 is amplified by an integrating amplifier, which generates an electrical pulse whose height is proportional to the number of emitted electrons. The height of the pulse is discriminated by a single channel analyzer (SCA). A logical AND is performed between the output of the SCA and the output of the photodiode to reject all detection events which occur outside of the pulse duration. Figure 2a shows a pulse height histogram of VLPC 1. This histogram features a series of peaks corresponding to different photon number states. The SCA can select pulse heights corresponding to one, two, three, and four photon events. The decision window set by the SCA for each photon number event is shown by the shaded areas in the figure.

VLPC 2 is placed in the idler arm, and used to verify that the correct photon number state was generated. The output of VLPC 2 is amplified, and connected to the signal input of a boxcar integrator. On each pulse from the SCA, the output of VLPC 2 is integrated over a 2$`\mu `$s window, which is sufficiently large to encompass the entire electrical pulse (determined by the bandwidth of subsequent amplifiers). Figure 2b shows a pulse area distribution of VLPC 2 with no post-selection from VLPC 1, at a pump excitation power of 40$`\mu `$W. This distribution also features a series of peaks corresponding to the different photon numbers, starting with the first peak which is a zero photon event. In order to calculate the photon number distribution, each peak is fit to a gaussian. The area of each gaussian is normalized by the total area of all the peaks. The calculated photon number distribution is shown in the inset.

Both VLPCs have imperfect detection efficiencies due to internal detection losses of the VLPC itself, as well as external losses from the collection optics. The efficiency is measured by comparing the coincidence rate between the two detectors to the singles count rates. From this measurement we determine the detection efficiency of VLPC 1 to be 0.68, and VLPC 2 to be 0.58.

The detection efficiency of the monitoring detector (VLPC 2) can be corrected for, because we are interested mainly in how many photons were actually present in the idler arm, not how many have been detected. To do this, we can define $`p_j`$ as the probability that j photons were generated in the idler pulse, and $`f_i`$ as the probability that i photons were detected. These two probabilities are related by the linear loss model $`f_i=_{j=i}^{\mathrm{}}\left(\genfrac{}{}{0pt}{}{j}{i}\right)\eta ^i\left(1\eta \right)^{ji}p_i`$. The above transformation needs to be inverted. In order to do this we truncate the photon number distribution at some number $`n`$, which is sufficiently large so that $`p(n+1)0`$ is a good approximation. The initial and final probability distributions are then related by a constant matrix, which can simply be inverted. This matrix can also be modified to account for dark counts in a straightforward way.

The detection efficiency of the triggering detector (VLPC 1) plays a more subtle role in the experiment. Detection losses can result in a higher order photon number state being misinterpreted by the detector as the correct photon number. Because of this, the probability distribution in the idler arm will no longer be an exact photon number, but a mixture of the desired number plus higher order number states. This results in a degradation of the fidelity, defined as the overlap between the desired state and the actual generated state. For a state containing $`n`$ photons the fidelity is simply the probability that $`n`$ photons are generated.

When the pumping intensity is low, the efficiency of the triggering detector does not play a major role. If the triggering detector sees $`n`$ photons, it is true with very high probability that the same number of photons are present in both arms. The probability that there were more photons present is extremely low, because low pumping intensity makes the probability of generating higher order photon pairs negligible. Figure 3 shows the experimental results at this low pumping intensity regime. The figure shows the photon number distributions measured by VLPC 2, after correcting for detection efficiency and dark counts, when VLPC 1 post-selects a one, two, three, and four photon event. For one, two, and three photon post-selection, a nearly ideal photon number state is generated. For four photon post-selection however, there are contributions from three and five photon number states. These contributions are attributed to the smearing between the four photon peak and its next nearest neighbors in the pulse height histogram of VLPC 1 (Fig. 2a). The smearing is caused by buildup of multiplication noise, which puts a limit on the photon number resolution at higher numbers Waks et al. . Note that, in a few cases, the corrected probability distribution becomes negative. This erroneous effect is caused by numerical errors in the probability reconstruction.

As we increase the pumping power, the imperfect detection efficiency of VLPC 1 will result in degraded fidelity, as was previously discussed. On the other hand, higher pump powers will increase the probability that the correct number of photons were generated, and hence the generation rate of the desired photon number state. This leads to a natural tradeoff between the state fidelity and generation rate of photon number states. The extent of the tradeoff is determined mainly by the detection efficiency of VLPC 1. In order to determine this tradeoff, we measure the fidelity and generation rate as a function of pump power for the four different photon number states. The results are shown in Figure 4. For all four cases, increasing the pump power results in higher generation rates but decreased fidelity. In the figure, four data points are labelled as a, b, c, and d. These points designate the pumping intensity before the fidelity begins to drop for one, two, three, and four photon number states. The generation rate and fidelity at these four points are: a - 11800 Hz and $`F_1`$=1.0, b -1100 Hz and $`F_2`$=0.94, c - 160 Hz and $`F_3`$=0.95, d - 84 Hz and $`F_4`$=0.6. The above generation rates are already large enough to be useful for practical experiments. These rates are achieved using only a 45kHz pump repetition rate. The rates could be significantly enhanced by simply using a pump laser with a faster repetition rate.

We conclude by discussing some of the future extensions of this work. There are two main drawbacks to the experimental scheme we present in this letter. First, the setup does not allow for generation of photon number states on demand. The time in which the number state is generated is known but uncontrollable. This problem can be solved in principle by setting up a large parallel array of photon number generators along with an optical switch. By making the array large, the probability that at least one of the number state generators will generate the right photon number can be made arbitrarily close to unity. Of course, with the current setup such a solution is impractical, but perhaps the engineering of VLPC arrays, along with more compact two photon sources based on waveguide technology, may make such a solution possible.

The second drawback of our experiment is that the pump pulse to the down-conversion crystal has a 20ns duration, as opposed to the coherence time of the down-converted photons which is on the order of femtoseconds. This means that, although the states we generate are eigenfunctions of the photon number operator, the quantum mechanical wavefunctions of the photons is not fixed. In fact the wavefunction of each photon will generally fluctuate from shot to shot. A more ideal source would generate a state in which both the number and wavefunction of the photons is constant. The best case scenario would be the generation of a so-called n-photon Fock state, which is a state containing n-photons with identical wavefunctions. Our setup could be extended to generate such a state by replacing the pumping laser with a femtosecond pulse laser, so that the pulse duration is on the order of the coherence length of the photons. This could open up the door for interesting new experiments in quantum information using photon number states.


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# The rôle of electronic correlations on the the phonon modes of MnO and NiO

## Abstract

The possibility of magnetic-order induced phonon anisotropy in single crystals of MnO and NiO is investigated using inelastic neutron scattering. Below $`\mathrm{T}_\mathrm{N}`$ both compounds exhibit a splitting in their transverse optical phonon spectra of approximately 10%. This behavior illustrates that, contrary to general assumption, the dynamic properties of MnO and NiO are substantially non-cubic.

The failure of *ab initio* approximations to correctly incorporate many-body effects such as electron exchange and correlation is an issue at the core of contemporary solid-state research. Although these effects are important in almost all solids, they are essential for a proper description of co-operative phenomena such as antiferromagnetism, charge-ordering, superconductivity and colossal magnetoresistance.

Despite the status of MnO and NiO as benchmark materials for the study of correlated electron systems and frequent use in first principles electronic structure investigations, many aspects of the physics of the 3$`d`$ transition metal monoxides requires better theoretical explanation. For example, the basic physical properties predicted for MnO and NiO (e.g. band-gaps, distortion angles and phonon spectra) differ radically depending upon which techniques have been employed M.D.Towler et al. (1994); S.Massidda et al. (1999).

In a recent publication comparing different *ab initio* and model band-structure models on MnO (including local spin-density approximation (LSDA) and ‘LSDA + model’ calculations), Massidda *et al.* S.Massidda et al. (1999) suggest that although the electronic density in MnO is approximately cubic, the lowering of symmetry associated with antiferromagnetic ordering can induce an electronic response that is significantly non-cubic. This results in the dynamical properties of the system, such as transverse optical (TO) phonons, exhibiting substantial “*magnetic-order induced anisotropy*”.

One of the predicted features from such calculations is that the zone center (ZC) optical phonon modes should split depending on their polarization. A higher energy occurs for a mode polarized along the  direction and a lower energy for degenerate modes polarized in the orthogonal ferromagnetic plane, see Table 1. Massidda *et al.* S.Massidda et al. (1999) find that this effect arises *solely due to the magnetic ordering*, and would even occur in the absence of any rhombohedral distortion. By comparing results from four different approximation schemes, lower and upper bounds of 3-10% for the magnitude of the splitting are estimated. Variations in the splitting due to the distortion are calculated to be less than 1%.

Although similar theoretical calculations have not yet been performed for NiO, the results of measurements of NiO TO phonons are predicted to be qualitatively similar to those in MnO. Indeed, it has been suggested that the splitting in NiO may be even greater in absolute magnitude than in MnO due to the larger magnetic super-exchange S.Massidda (2001).

MnO and NiO are classic examples of type-II antiferromagnets possessing the cubic rock-salt structure. Below $`\mathrm{T}_\mathrm{N}`$, exchange-striction causes contraction of the unit cell along the  direction perpendicular to the ferromagnetic planes S.Greenwald and J.S.Smart (1950); W.L.Roth (1958). This results in a small deviation from the perfect cube of 0.62o for MnO and 0.1o for NiO. In MnO the magnitude and temperature dependence of the rhombohedral distortion has been well characterized B.Morosin (1970). Since the distortion is considered weak, investigations of the dynamic properties of MnO and NiO are frequently carried out under the assumption of perfect cubic symmetry.

To examine the possibility of anisotropy in the optical phonon spectra of MnO (and NiO), (and distinguish these effects from any changes due to the small rhombohedral distortion), TO phonons in single crystals of MnO and NiO were measured using inelastic neutron scattering. MnO and NiO crystals used in these experiments were grown using the floating-zone method at a rate of 10 mm/h in 1 atm of Argon. Both crystals consisted of small cylindrical pillars approximately 0.5 cm in diameter and 2 cm long.

The optic phonons, magnon modes and sublattice magnetization of our MnO crystal were studied using the IN8 triple axis spectrometer located at the high-flux reactor at the Institut Laue Langevin, (ILL), Grenoble. The Cu(111) reflection was used as monochromator in combination with a pyrolytic graphite (PG) filter, PG(002) analyzer and Soller collimation. Measurements were made using constant-*Q* scans, with a fixed final energy of $`\mathrm{E}_\mathrm{f}`$ = 14.69 meV, ($`\mathrm{k}_\mathrm{f}`$ = 2.66 Å). The crystal was oriented with the  direction vertical.

Operating the spectrometer in diffraction mode, measurements of the temperature evolution of the magnetic ($`\frac{3}{2}`$ $`\frac{3}{2}`$ $`\frac{3}{2}`$) reflection were obtained. The observed value of $`\mathrm{T}_\mathrm{N}`$= 118 K, measured during slow cooling at a rate of 1o every 3 mins, was in excellent agreement with previous reports D.Bloch et al. (1974); R.Boire and M.F.Collins (1977); B.F.Woodfield et al. (1999).

Compression along equivalent  directions in the bulk crystal creates four possible orientations of twinning, or ‘T-domain’, each with a distinct 3*D* orientation of magnetic Brillouin zone. The scattering plane of the sample is then comprised of a superposition of slices through each of the magnetic Brillouin zones M.Kohgi et al. (1974); M.Bonfante et al. (1972). The T-domains in our MnO sample were characterized using spin wave dispersion measurements obtained at 4.3 K. Proceeding by continuity from unambiguously defined reciprocal lattice positions, the full magnon dispersion curves of our multi-T-domain MnO crystal along the ,  and  directions were determined. From these dispersion curves it was possible to deduce that the crystal was comprised of two approximately equally populated twinned domains.

Measurements of transverse optical (TO) phonon modes were obtained along several directions in the crystal (,  and ) and around various Brillouin zone (BZ) centers (Q=(004), Q=(113), Q=(222) and Q=(220)). At 4.3 K, splitting of the modes was consistently observed in all of the measured BZ’s. Representative zone center scans obtained at 4.3 K and 300 K are shown in figure 1.

The energy of the 300 K ZC excitation occurs at 33.2 $`\pm `$ 0.2 meV with a full-width half-maximum of 2.87 $`\pm `$ 0.5 meV. In the antiferromagnetically ordered phase at 4.3 K two modes were consistently observed at energies of E = 33.3 $`\pm `$ 0.2 meV and E = 36.4 $`\pm `$ 0.3 meV. The zone center splitting at 4.3 K is approximately 3.1 meV, which corresponds to 9.3% of the energy of the original (degenerate) mode above $`\mathrm{T}_\mathrm{N}`$.

The uncertainty in reduced wave vector, q, arising from the rhombohedral distortion was estimated for the measured Brillouin zones to range between *d*q= 0.02 and *d*q= 0.04 r.l.u. At the ZC position where the dispersion of the TO modes is virtually flat the effects of an uncertainty of this magnitude on the energy spectra would be negligible. We therefore conclude that the splitting is a real effect which can not be attributed to the small static structural distortion.

The low temperature data were fitted to a pair of Gaussian peaks of identical width. Peaks at each of the zone-centers possessed widths of 2.43 $`\pm `$ 0.6 meV. The ratio of the integrated intensities of the two peaks was approximately 2:1 in favor of the lower mode at each of the measured zone-centers. The total integrated intensity of the two (split) peaks and the initial single peak were approximately conserved on passing through TN. This behavior is consistent with the creation of a new higher energy mode arising from excitations polarized along the  direction S.Massidda et al. (1999). Although both modes are observed at a slightly higher energy than the ‘LSDA + model’ value (table 1), the observed 9% anisotropy is within the suggested bounds of 3-10% S.Massidda et al. (1999).

Plots of the temperature evolution of the ZC phonon anisotropy are shown in figures 1(inset) and 2. The splitting in the phonon spectra is clearly associated with the transition to magnetic order. Asymmetry in peaks measured at 130 K and 150 K (fig. 2) hints that a small amount of anisotropy may be present at temperatures substantially above $`\mathrm{T}_\mathrm{N}`$, (i.e. when the unit cell is thought to be perfectly cubic B.Morosin (1970)).

Dispersion curves obtained for TO modes along the  and  directions in MnO at 4.3 K and 150 K are presented in figures 3 and 4(a). Although the splitting extends throughout the Brillouin zone, the distortion is most clearly resolved at values of reduced wavevector (q) close to the zone-center, see figure 3(a). Further from the ZC (where the scattering is less intense and the spectrometer may be less focussed), the presence of two modes is revealed by the deviation of the shape of the overall peak from Gaussian, and its increased width in comparison to the same measurement at 150 K, fig. 3(b). As the crystal consisted of two domains, intensity from differently polarized magnons could be seen reasonably well in all of the measured directions.

The notion that a small splitting may be present in MnO above TN is supported by the dispersion measured at 150 K, shown in figure 4(a). The 150 K data (open circles) was fitted to a single Gaussian peak but found to be slightly raised in energy in comparison to the lower energy mode (filled circles), and to the room temperature data of Wagner *et al.* V.Wagner et al. (1976) (square markers).

To investigate whether a similar phonon splitting could be observed in the room temperature dispersion of NiO, a preliminary survey was performed using the HB1 spectrometer located at the High Flux Isotope Reactor, Oak Ridge National Laboratory, USA. As before, the crystal was oriented with the  direction vertical, and the spectrometer operated in constant-*Q* mode with fixed final energy $`\mathrm{E}_\mathrm{f}`$ = 14.73 meV.

Our NiO crystal was of less reliable quality than the MnO crystal, and spurious ‘Bragg-Bragg-inelastic’ scattering was observed. This limited our measurements to regions surrounding ZC’s located at Q=(220) and Q=(222). Room temperature TO dispersion curves for the  and  directions of antiferromagnetically ordered NiO are presented in figure 4(b)(filled round markers). Squares represent the data of Reichardt *et al.* W.Reichardt et al. (1975), also collected at room temperature.

In our NiO data two peaks were observed occurring at energies of around 45 and 50 meV. This is in reasonable agreement with previous reports indicating a single TO excitation with an energy varying between 46 and 50 meV R.A.Coy et al. (1976); W.Reichardt et al. (1975). This ZC splitting of approximately 5 meV is larger in absolute magnitude than observed in MnO, (despite the smaller distortion angle), and corresponds to roughly 10% of the average energy of the two modes.

Since both this data and data of previous researchers were collected at room temperature (where NiO is already magnetically ordered) it is necessary to ask why this effect had not been been observed before. We suggest that the cause of the discrepancy lies in the poor energy resolution and flux of the early spectrometers, particularly at high energy transfers. These experimental factors coupled with the expectation of a single (degenerate) phonon mode may have lead to fitting of the data to a single peak.

In summary, we have shown that below $`\mathrm{T}_\mathrm{N}`$ the dynamic properties of MnO and NiO are substantially non-cubic. The magnitude of the zone center anisotropy observed in transverse optic modes in MnO and NiO corresponds to 9-10% of the energy of the original (degenerate) phonon. The intensities of the two optic modes in MnO and the form of the dispersion are consistent with the picture of the splitting described by recent band structure calculations where a higher energy mode is predicted for phonons polarized along  S.Massidda et al. (1999).

In MnO, we have shown that the splitting of the phonon modes is associated with the transition to magnetic order. Although exchange-striction induces a small rhombohedral distortion, this static distortion in itself is too weak to explain the magnitude of the anisotropy. Furthermore, NiO possesses a significantly smaller distortion angle than MnO but a larger (5 meV) splitting. The strain associated with exchange-striction (even in the undistorted structure considered by Massidda *et al*. S.Massidda et al. (1999)) defines a new potential well for the atoms, which may significantly influence the dynamics. However, further research will be required to investigate the mechanism underlying the splitting and determine how the magnetic-order induced phonon anisotropy in MnO and NiO is mediated. Theoretical studies investigating the rôle of electronic correlations on the phonon modes of NiO are particularly encouraged. Although theorists have not provided any calculations for NiO as yet, we are confident that our preliminary findings of an even larger splitting will provide a new impetus for them to do so.

We thank Sandro Massidda, Julie Staunton, Jim Hague, Randy Fishman and Oleg Petrenko for useful discussions, and the Engineering and Physical Sciences Research Council, UK, for financial support.