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# Independent sets in certain classes of (almost) regular graphs

## 1. Introduction

Let $`I(G)`$ denote the number of independent sets of a graph $`G`$. This number can be determined for some special classes of graphs (see  for a survey). For instance, $`I(G)`$ was studied for *grid graphs* (see ), *multipartite complete graphs*, and *path-* and *cyclic-schemes*. Many of these numbers are given by certain combinations of Fibonacci numbers, some others by Lucas numbers.

In this paper, we study four classes of graphs. To define these classes, we recall that a *line graph* $`L(G)`$ of a graph $`G`$ is obtained by associating a vertex with each edge $`G`$ and connecting two vertices with an edge if and only if the corresponding edges of $`G`$ are adjacent. Also, recall that a *cycle graph* $`C_{\mathrm{}}`$ is a graph on $`\mathrm{}`$ nodes containing a single cycle through all nodes.

We now give the definitions of our classes.

Let $`G_{\mathrm{}}^1=C_{\mathrm{}}`$, the $`\mathrm{}`$-cycle graph. We obtain $`G_{\mathrm{}}^2`$ by superimposing the line graph $`L(G_{\mathrm{}}^1)`$ onto the graph $`G_{\mathrm{}}^1`$, that is splitting each edge of $`G_{\mathrm{}}^1`$ with the corresponding vertex of $`L(G_{\mathrm{}}^1)`$ and then adding the edges of $`L(G_{\mathrm{}}^1)`$. More generally, $`G_{\mathrm{}}^n`$ is obtained by superimposing $`L^{n1}(G_{\mathrm{}}^1)`$ onto $`G_{\mathrm{}}^{n1}`$. For example, in Figure 1, we have the graphs $`G_3^4`$ and $`G_4^3`$, respectively, if one ignores the dashed edges. Clearly, all but $`n`$ of the $`\mathrm{}^n`$ nodes of $`G_{\mathrm{}}^n`$ have degree 4, and we say that $`G_{\mathrm{}}^n`$ is an *almost* 4-regular graph.

The graph $`R_{\mathrm{}}^n`$ is obtained from $`G_{\mathrm{}}^n`$ by duplicating the edges of $`G_{\mathrm{}}^1`$. For example, in Figure 1, the extra edges are the dashed edges, and by adding them we get the graphs $`R_3^4`$ and $`R_4^3`$ respectively. So, to get $`R_{\mathrm{}}^n`$ we add $`\mathrm{}`$ additional edges to $`G_{\mathrm{}}^n`$, and it is easy to see that $`R_{\mathrm{}}^n`$ is a 4-regular graph.

Let $`K_{\mathrm{}}^1=K_{\mathrm{}}`$, a complete graph on $`\mathrm{}`$ nodes. Put the $`\mathrm{}`$ nodes of $`K_{\mathrm{}}^1`$ on a circle and draw the remaining $`(\mathrm{}1)!\mathrm{}`$ edges. Call the first $`\mathrm{}`$ edges *external* and the remaining edges, *internal*. Then construction of $`K_{\mathrm{}}^n`$ is similar to that of $`G_{\mathrm{}}^n`$, except that:

1. The basis of the construction is now $`K_{\mathrm{}}`$, rather than $`C_{\mathrm{}}`$.
2. On each iteration $`i`$, the graph superimposed onto $`R_{\mathrm{}}^{i1}`$ is not $`L^i(K_{\mathrm{}})`$ but rather the complete graph on the nodes of $`L^i(C_{\mathrm{}})`$, the line graph of $`C_{\mathrm{}}`$ formed by the external edges of $`K_{\mathrm{}}^1`$.

In Figure 2, we show how to construct $`K_4^3`$ from $`K_4^2`$ (the dashed edges should be ignored). In that figure, the external edges of respective complete graphs are in bold. We also remark that the internal edges do not intersect each other. Moreover, Figure 2 suggests a convenient way of representing $`K_4^3`$, where each node of the graph lies only on internal edges incident with that node. We achieve that by putting the nodes of the line graphs under consideration off the centers of the corresponding edges. Indeed, if we were using the centers of the external edges, $`K_4^3`$ would look as in Figure 3, which is misleading since, for example, the node $`a`$ does not lie on the edge $`bc`$.

The graph $`K_{\mathrm{}}^n`$ is almost $`(\mathrm{}+1)`$-regular, since all but $`\mathrm{}`$ of its $`\mathrm{}^n`$ nodes have degree $`\mathrm{}+1`$. Moreover, it follows from our definitions that $`K_3^n=G_3^n`$.

The graph $`P_{\mathrm{}}^n`$ is obtained from the graph $`K_{\mathrm{}}^n`$ by duplicating the external edges in the graph $`K_{\mathrm{}}^1`$. For example, in Figure 2, the extra edges are the dashed edges, and by adding them we construct the graph $`P_4^3`$ from $`K_4^3`$. So to get $`P_{\mathrm{}}^n`$ we add $`\mathrm{}`$ extra edges to $`K_{\mathrm{}}^n`$, and it is easy to see that $`P_{\mathrm{}}^n`$ is an $`(\mathrm{}+1)`$-regular graph.

Let $`g_{\mathrm{}}(n)`$, $`r_{\mathrm{}}(n)`$, $`k_{\mathrm{}}(n)`$ and $`p_{\mathrm{}}(n)`$ denote the number of independent sets in the graphs $`G_{\mathrm{}}^n`$, $`R_{\mathrm{}}^n`$, $`K_{\mathrm{}}^n`$ and $`P_{\mathrm{}}^n`$ respectively. In our paper we study these numbers. In Section 3, we give an algorithm for calculating all these numbers. For the numbers $`p_{\mathrm{}}(n)`$, we provide an explicit generating function (see Theorem 3.4). However, in order to illustrate our approach to the problem, in Section 2 we consider $`g_3(n)`$ and find an explicit formula for it.

Our choice of the graphs to study was motivated by the so called de Bruijn graphs, which are defined as follows. A de Bruijn graph is a directed graph $`\stackrel{}{G}_n=\stackrel{}{G}_n(V,E)`$, where the set of vertices $`V`$ is the set of all the words of length $`n`$ in a finite alphabet $`A`$, and there is an arc from $`v_i=(v_{i1},\mathrm{},v_{in})`$ to $`v_j=(v_{j1},\mathrm{},v_{jn})`$ if

$`v_{i2}=v_{j1}`$, $`v_{i3}=v_{j2},\mathrm{},v_{in}=v_{j(n1)}`$,

that is when the words $`v_i`$ and $`v_j`$ overlap by $`(n1)`$ letters.

The de Bruijn graphs were first introduced (for the alphabet $`A=\{0,1\}`$) by de Bruijn in 1944 for enumerating the number of code cycles. However, these graphs proved to be a useful tool for various problems related to the subject of combinatorics on words (e.g. see ). It is known that the graph $`\stackrel{}{G}_n`$ can be defined recursively as $`\stackrel{}{G}_n=L(\stackrel{}{G}_{n1})`$. The authors were interested in studying other graphs defined recursively using the operation of taking line graphs (with natural bases), which could give interesting applications. Also, with our choice of graphs ($`G_{\mathrm{}}^n`$ and $`K_{\mathrm{}}^n`$), it is natural to complete them to regular graphs ($`R_{\mathrm{}}^n`$ and $`P_{\mathrm{}}^n`$) and study these graphs. It turns out that there are combinatorial interpretations (relations to other combinatorial objects) for the number of independent sets for some of our graphs, and we mention these relations in Sections 2 and 3. Moreover, we construct a direct bijection describing such a relation for $`P_4^n`$ (see Proposition 3.5 and the discussion that follows).

## 2. The numbers $`g_3(n)`$.

Let us first find an explicit formula for $`g_3(n)`$.

It is clear that for any independent set of the graph $`G_3^n`$, we can label a node of $`G_3^n`$ 1 if this node is in the independent set, and label it 0 otherwise. Thus, our purpose is to count the number of triangles having either 0 or 1 in each node and such that no two adjacent nodes are both assigned 1s. We call such triangles *legal*.

In order to get a recursion for $`g_3(n)`$, we introduce three auxiliary parameters $`a_n`$, $`b_n`$, and $`c_n`$, which are the numbers of legal triangles that, up to rotation, have specific numbers in the nodes of the biggest triangle (see Figure 4). Since we consider only legal graphs, the 1s in the nodes of the biggest triangle induces 0s in certain nodes of a smaller triangle (this 0s are shown in Figure 4).

Considering all the possibilities for the numbers of the biggest triangle, we have that

$$g_3(n)=g_3(n1)+3a_n+3b_n+c_n,$$

where $`g_3(n1)`$ corresponds to all 0s, and we have multiple 3 two times because of possible rotations. Similarly, we get that $`a_n=g_3(n2)+a_{n1}`$, and $`b_n=c_n=g_3(n2)`$. This leads to the recursion

(1)
$$g_3(n)=2g_3(n1)+6g_3(n2)4g_3(n3),$$

which, under the same initial conditions, is equivalent to the recursion

(2)
$$g_3(n)=4g_3(n1)2g_3(n2).$$

We define $`g_3(0)=1`$, since we associate the graph $`G_3^0`$ with the empty graph, in which case there is only one independent set, the empty set. Thus,

$$g_3(n)=\frac{1}{2\sqrt{2}}\left((2+\sqrt{2})^{n+1}(2\sqrt{2})^{n+1}\right),$$

and the generating function for the numbers $`g_3(n)`$ is $`1/(14x+2x^2)`$. The initial values for the numbers $`g_3(n)`$ are:

$$1,4,14,48,164,560,1912,6528,22288,76096,\mathrm{}.$$

Preceded by 0, the sequence $`\{g_3(n)\}`$ is the binomial transform of the Pell numbers

$$P_n=\frac{(1+\sqrt{2})^n(1\sqrt{2})^n}{2\sqrt{2}}$$

(see \[10, A007070\]). These numbers can also be interpreted as maximum bets in a poker game (also see \[10, A007070\]), where the first player bets 1 dollar into a pot and the $`i`$th player bets the amount of the $`(i1)`$st player’s bet plus the resulting amount of money in the pot. Then the number of dollars $`d_n`$ in the pot after $`n`$ bets is given by

$$d_n=2(d_{n1}+(d_{n1}d_{n2}))=4d_{n1}2d_{n2},d_0=1,d_1=4,$$

which yields $`d_n=g_3(n)`$.

We remark that it would be interesting to obtain recurrence (2) directly from the graph $`G_3^n`$, rather than via recurrence (1). Unfortunately, we were unable to do this.

## 3. An algorithm for calculating $`g_{\mathrm{}}(n)`$, $`r_{\mathrm{}}(n)`$, $`k_{\mathrm{}}(n)`$ and $`p_{\mathrm{}}(n)`$

In this section we present an algorithm for calculating $`g_{\mathrm{}}(n)`$, $`r_{\mathrm{}}(n)`$, $`k_{\mathrm{}}(n)`$ and $`p_{\mathrm{}}(n)`$ by using the transfer matrix method (see \[9, Theorem 4.7.2\]).

### 3.1. An algorithm for calculating $`g_{\mathrm{}}(n)`$

In this section, we use the transfer matrix method to obtain an information about the sequence of $`g_{\mathrm{}}(n)`$.

Similarly to Section 2, for any independent set of the graph $`G_{\mathrm{}}^n`$, consider a labeling of $`G_{\mathrm{}}^n`$, where the nodes of of the independent set are labeled $`1`$ and the remaining nodes are labeled 0. For a given graph $`G_{\mathrm{}}^n`$, we define the *$`n`$-th level* of $`G_{\mathrm{}}^n`$ to be $`G_{\mathrm{}}^nG_{\mathrm{}}^{n1}`$, which is isomorphic to $`G_{\mathrm{}}^1`$. Thus, we may think of an independent set of the graph $`G_{\mathrm{}}^n`$ as assembled from elements chosen on each level, making sure that when we add a new level, we create no conflict with the previous level.

The collection $`_{\mathrm{}}`$ of possible level labelings is the set of all $`(0,1)`$ $`\mathrm{}`$-vectors $`𝐯=(v_1,\mathrm{},v_{\mathrm{}})`$. It will be convenient to define $`v_{\mathrm{}+1}:=v_1`$. Then $`𝐯=(v_1,\mathrm{},v_{\mathrm{}})`$ and $`𝐰=(w_1,\mathrm{},w_{\mathrm{}})`$ in $`_{\mathrm{}}`$ are a possible consecutive pair of levels in an independent set of $`G_{\mathrm{}}^n`$ (with $`𝐰`$ following $`𝐯`$) if and only if

(3)
$$v_i=v_{i+1}=0\text{ or }w_i=0,\text{ where }i=1,2,\mathrm{},\mathrm{}.$$

Thus, to obtain any independent set in the graph $`G_{\mathrm{}}^n`$, we begin with a vector of $`_{\mathrm{}}`$, then keep adjoining each next vector $`𝐰_{\mathrm{}}`$ so that it satisfies (3) together with the previously chosen vector $`𝐯_{\mathrm{}}`$, until $`n`$ vectors have been selected.

We define a matrix $`G=G_{\mathrm{}}`$, the transfer matrix of the problem, as follows. $`G`$ is a $`2^{\mathrm{}}\times 2^{\mathrm{}}`$ matrix of $`0`$s and $`1`$s whose rows and columns are indexed by vectors of $`_{\mathrm{}}`$. The entry of $`G`$ in position $`(𝐯,𝐰)`$ is $`1`$ if the ordered pair of vectors $`(𝐯,𝐰)`$ satisfies (3), and is $`0`$ otherwise. $`G`$ depends only on $`\mathrm{}`$, not on $`n`$. Hence, the number of independent sets of $`G_{\mathrm{}}^n`$, $`g_{\mathrm{}}(n)`$, is the first entry of the vector $`G^n(u_1,\mathrm{},u_2^{\mathrm{}})^T`$, where $`u_i=1`$ if there are the $`i`$th vector $`𝐯`$ in the collection $`_{\mathrm{}}`$ has no two consecutive $`1`$s, even after wrapping, (i.e. $`v_i+v_{i+1}1`$ for all $`i=1,2,\mathrm{},\mathrm{}`$), and $`u_i=0`$ otherwise. Hence,

$$g_{\mathrm{}}(n)=(1,0,\mathrm{},0)G^n(u_1,\mathrm{},u_2^{\mathrm{}})^T.$$

For instance, when $`\mathrm{}=3`$, the possible level vectors are

$$(0,0,0),(0,0,1),(0,1,0),(0,1,1),(1,0,0),(1,0,1),(1,1,0),(1,1,1),$$

except for the last level, where we only have

$$(0,0,0),(0,0,1),(0,1,0),(1,0,0).$$

If we index the rows and the columns of the transfer matrix $`G`$ in this order, then we get

$$G=\left[\begin{array}{cccccccc}1& 1& 1& 1& 1& 1& 1& 1\\ 1& 0& 0& 0& 1& 0& 0& 0\\ 1& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 1& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right].$$

The vector $`(u_1,\mathrm{},u_2^{\mathrm{}})^T`$ in this case is $`(1,1,1,0,1,0,0,0)`$. If we now find the first entry of the vector $`(IxG)^1(1,1,1,0,1,0,0,0)^T`$, where $`I`$ is the unit matrix, then we get that the generating function for $`g_3(n)`$ is given by $`1/(14x+2x^2)`$. We obtain the results for larger $`\mathrm{}`$ similarly.

###### Theorem 3.1.

The generating functions for the numbers $`g_4(n)`$, $`g_5(n)`$ and $`g_6(n)`$ are given, respectively, by

$$\frac{1+4xx^22x^3}{13x14x^2+15x^3+7x^4},$$

$$\frac{(1+x)(1+5x8x^2)}{15x30x^2+69x^3+31x^422x^5},$$

$$\frac{1+10x12x^250x^3+10x^4+20x^512x^6}{18x66x^2+280x^3+178x^4532x^584x^6+108x^7}.$$

We remark that the algorithm for finding the generating function for $`g_{\mathrm{}}(n)`$ has been implemented in Maple, and yielded explicit results for $`\mathrm{}6`$.

### 3.2. An algorithm for calculating $`r_{\mathrm{}}(n)`$

In this section we use the transfer matrix method to obtain an information about the numbers $`r_{\mathrm{}}(n)`$. This case is similar to that of $`g_{\mathrm{}}(n)`$ with some small differences.

We partition $`R_{\mathrm{}}^n`$ into levels just as in the case of $`G_{\mathrm{}}^n`$, so the $`n`$-th level of $`R_{\mathrm{}}^n`$ is $`R_{\mathrm{}}^n\backslash R_{\mathrm{}}^{n1}`$.

The collection of possible levels $`_{\mathrm{}}`$ is the set of all $`\mathrm{}`$-vectors $`𝐯`$ of $`0`$s and $`1`$s such that there no consecutive $`1`$s in $`𝐯`$, that is, $`v_i+v_{i+1}2`$ (where we define $`v_{\mathrm{}+1}:=v_1`$). Clearly, the set $`_{\mathrm{}}`$ contains exactly $`L_{\mathrm{}}`$ vectors where $`L_{\mathrm{}}`$ is the $`\mathrm{}`$th Lucas number. For instance, $`_3`$ contains the vectors $`(0,0,0)`$, $`(0,0,1)`$, $`(0,1,0)`$, and $`(1,0,0)`$.

The condition that vectors $`𝐯`$ and $`𝐰`$ in $`_{\mathrm{}}`$ are a possible consecutive pair of levels in an independent set of $`R_{\mathrm{}}^n`$ is given by (3) just as for $`G_{\mathrm{}}^n`$. To obtain any independent set in the graph $`R_{\mathrm{}}^n`$, we begin with a vector of $`_{\mathrm{}}`$, then keep adjoining each next vector $`𝐰_{\mathrm{}}`$ so that it satisfies (3) together with the previously chosen vector $`𝐯_{\mathrm{}}`$, until $`n`$ vectors have been selected.

We define a matrix the transfer matrix of the problem $`R=R_{\mathrm{}}`$ in the same way as $`G=G_{\mathrm{}}`$ in subsection 3.1. Then $`R`$ is an $`L_{\mathrm{}}\times L_{\mathrm{}}`$ matrix, and the number of independent sets of $`R_{\mathrm{}}^n`$, $`r_{\mathrm{}}(n)`$, is the first entry of of the vector $`R^n\mathrm{𝟏}`$ where $`\mathrm{𝟏}=(1,1,\mathrm{},1)`$. Hence,

$$r_{\mathrm{}}(n)=(1,0,\mathrm{},0)R^n\mathrm{𝟏}.$$

For instance, when $`\mathrm{}=3`$, the possible level vectors in an independent set are

$$(0,0,0),(0,0,1),(0,1,0),(1,0,0).$$

If we indexed the rows and columns in this order, then the transfer matrix is

$$R=\left[\begin{array}{cccc}1& 1& 1& 1\\ 1& 0& 0& 1\\ 1& 1& 0& 0\\ 1& 0& 1& 0\end{array}\right].$$

If we find the first entry of the vector $`(IxG)^1\mathrm{𝟏}`$, we get that the generating function for $`r_3(n)`$ is given by $`\frac{1+2x}{12x2x^2}`$. The initial values for the numbers $`r_3(n)`$ are

$$1,4,10,28,76,208,568,1552,4240,\mathrm{}$$

This sequence appears as A026150 in .

Similarly to the case $`\mathrm{}=3`$, we obtain the following results for $`\mathrm{}=4,5,6`$.

###### Theorem 3.2.

The generating functions for the numbers $`r_4(n)`$, $`r_5(n)`$ and $`r_6(n)`$ are given, respectively, by

$$\frac{1+4x4x^2}{13x4x^2+4x^3},$$

$$\frac{1+7x6x^2}{14x8x^2+6x^3},$$

$$\frac{1+12x24x^2+8x^4}{(18x+4x^2+4x^3)(1+2x2x^2)}.$$

We remark that the algorithm for finding the generating function for $`r_{\mathrm{}}(n)`$ has been implemented in Maple, and yielded explicit results for $`\mathrm{}6`$.

### 3.3. An algorithm for calculating $`k_{\mathrm{}}(n)`$

In this section we use the transfer matrix method yet again to obtain information about the sequences $`k_{\mathrm{}}(n)`$. This case is also similar to that of $`g_{\mathrm{}}(n)`$, so we will only sketch it briefly.

We partition $`K_{\mathrm{}}^n`$ into levels just as in the case of $`G_{\mathrm{}}^n`$, so the $`n`$-th level of $`K_{\mathrm{}}^n`$ is $`K_{\mathrm{}}^n\backslash K_{\mathrm{}}^{n1}`$. The collection of possible levels $`_{\mathrm{}}`$ is the set of all $`\mathrm{}`$-vectors $`𝐯`$ of $`0`$s and $`1`$s. Vectors $`𝐯`$ and $`𝐰`$ in $`_{\mathrm{}}`$ are a possible consecutive pair of levels in an independent set of $`K_{\mathrm{}}^n`$ if they satisfy (3).

We define the transfer matrix of the problem, $`K=K_{\mathrm{}}`$, in the same way as $`G_{\mathrm{}}`$. Then $`K`$ is a $`2^{\mathrm{}}\times 2^{\mathrm{}}`$ matrix, and the number of independent sets of $`K_{\mathrm{}}^n`$, $`k_{\mathrm{}}(n)`$, is the first entry of of the vector $`K^n(u_1,\mathrm{},u_2^{\mathrm{}})^T`$ where $`u_i=1`$ if the $`i`$th vector in the collection $`_{\mathrm{}}`$ contains at most one nonzero entry. Hence,

$$k_{\mathrm{}}(n)=(1,0,\mathrm{},0)K^n(u_1,\mathrm{},u_2^{\mathrm{}})^T.$$

For instance, when $`\mathrm{}=3`$, the possible level vectors in an independent set are

$$(0,0,0),(0,0,1),(0,1,0),(0,1,1),(1,0,0),(1,0,1),(1,1,0),(1,1,1).$$

If we index the rows and columns in this order, then the transfer matrix is

$$K=\left[\begin{array}{cccccccc}1& 1& 1& 1& 1& 1& 1& 1\\ 1& 0& 0& 0& 1& 0& 0& 0\\ 1& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 1& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right].$$

If we find the first entry of the vector $`(IxK)^1(u_1,\mathrm{},u_2^{\mathrm{}})^T`$, we get that the generating function for $`k_3(n)`$ is given by $`\frac{1}{14x+2x^2}`$. In particular, we get that $`k_3(n)=g_3(n)`$ which is can also be seen directly from the definitions. Similarly, we have the following result.

###### Theorem 3.3.

The generating functions for the numbers $`k_4(n)`$, $`k_5(n)`$ and $`k_6(n)`$ are given, respectively, by

$$\frac{1+2x+3x^2}{13x14x^2+15x^3+7x^4},$$

$$\frac{1+x+12x^28x^3}{15x30x^2+69x^3+31x^422x^5},$$

$$\frac{1x+38x^272x^38x^4+30x^5}{18x66x^2+280x^3+178x^4532x^584x^6+108x^7}.$$

We remark that the algorithm for finding the generating function for $`k_{\mathrm{}}(n)`$ has been implemented in Maple, and yielded explicit results for $`\mathrm{}6`$.

### 3.4. An algorithm for calculating $`p_{\mathrm{}}(n)`$

In this section we use the transfer matrix method to obtain information about the sequences $`p_{\mathrm{}}(n)`$.

In this section we use the transfer matrix method once more to obtain information about the sequences $`k_{\mathrm{}}(n)`$. This case is also similar to that of $`g_{\mathrm{}}(n)`$, so we will only sketch it briefly.

We partition $`P_{\mathrm{}}^n`$ into levels the same way as $`G_{\mathrm{}}^n`$, so the $`n`$-th level of $`P_{\mathrm{}}^n`$ is $`P_{\mathrm{}}^n\backslash P_{\mathrm{}}^{n1}`$. The collection of possible levels $`_{\mathrm{}}`$ is the set of all $`\mathrm{}`$-vectors $`𝐯`$ of $`0`$s and $`1`$s. Vectors $`𝐯`$ and $`𝐰`$ in $`_{\mathrm{}}`$ are a possible consecutive pair of levels in an independent set of $`P_{\mathrm{}}^n`$ if they satisfy (3).

The collection of possible levels $`_{\mathrm{}}`$ is the set of all $`\mathrm{}`$-vectors $`𝐯=(v_1,\mathrm{},v_{\mathrm{}})`$ of $`0`$s and $`1`$s such that $`v_1+\mathrm{}+v_{\mathrm{}}1`$. Clearly, the set $`_{\mathrm{}}`$ contains exactly $`\mathrm{}+1`$ vectors which are $`(0,\mathrm{},0)`$ and $`(0,\mathrm{},0,1,0,\mathrm{},0)`$. The condition that vectors $`𝐯`$ and $`𝐰`$ in $`_{\mathrm{}}`$ are a possible consecutive pair of levels in an independent set of $`P_{\mathrm{}}^n`$ is given by (3).

We define the transfer matrix of the problem, $`P=P_{\mathrm{}}`$, in the same way as $`G_{\mathrm{}}`$, $`R_{\mathrm{}}`$ and $`K_{\mathrm{}}`$. We define a matrix $`P=P_{\mathrm{}}`$, the transfer matrix of the problem, as follows. $`P`$ is an $`(\mathrm{}+1)\times (\mathrm{}+1)`$ matrix of $`0`$s and $`1`$s whose rows and columns indexed by vectors of $`_{\mathrm{}}`$. Therefore, it is easy to see that

$$P=\left[\begin{array}{ccccccccc}1& 1& 1& 1& 1& \mathrm{}& 1& 1& 1\\ 1& 0& 0& 1& 1& \mathrm{}& 1& 1& 1\\ 1& 1& 0& 0& 1& \mathrm{}& 1& 1& 1\\ \mathrm{}& & & & & \mathrm{}& & & \mathrm{}\\ 1& 1& 1& 1& 1& \mathrm{}& 1& 0& 0\\ 1& 0& 1& 1& 1& \mathrm{}& 1& 1& 0\end{array}\right]_{(\mathrm{}+1)\times (\mathrm{}+1)}.$$

###### Theorem 3.4.

The generating function for $`p_{\mathrm{}}(n)`$ is given by

$$\underset{n0}{}p_{\mathrm{}}(n)x^n=\frac{1+2x}{1(\mathrm{}1)x2x^2}.$$

###### Proof.

We want to find the first entry of the vector $`(IxP)^1\mathrm{𝟏}`$ which means we must find the first row, say $`(e_1,\mathrm{},e_{\mathrm{}+1})`$, of the matrix $`(IxP)^1`$. By solving the system of equations

$$(IxP)^1(e_1,\mathrm{},e_{\mathrm{}+1})^T=(1,0,\mathrm{},0)^T,$$

we get that

$$e_1=\frac{1(\mathrm{}2)x}{1(\mathrm{}1)x2x^2}\text{ and }e_j=\frac{x}{1(\mathrm{}1)x2x^2}\text{ for }j2.$$

Hence, the first entry of the vector $`(IxP)^1\mathrm{𝟏}`$ is given by

$$\underset{i=1}{\overset{\mathrm{}+1}{}}e_i=\frac{1+2x}{1(\mathrm{}1)x2x^2}.$$

For instance, when $`\mathrm{}=3`$, the generating function for $`p_3(n)`$ is given by $`\frac{1+2x}{12x2x^2}`$. One can see, in particular, that $`p_3(n)=r_3(n)`$, which follows directly from the definitions.

In the case $`\mathrm{}=4`$ the initial values of the numbers $`p_4(n)`$ are

$$1,5,17,61,217,773,2753,9805,34921,124373,\mathrm{}$$

The same sequence turns out to appear in  (see \[10, A007483\]). Thus, the following proposition is true:

###### Proposition 3.5.

The number of independent sets in the graph $`P_4^n`$ is equal to the number of (possibly empty) subsequences of the sequence $`\{1,2,\mathrm{},2n+1\}`$ in which each odd member has an even neighbor.

Here, the neighbors of an integer $`m`$ are $`m1`$ and $`m+1`$. In the case $`n=1`$, the sequences appearing in the proposition are $`ϵ,2,23,12,123`$, where $`ϵ`$ is the empty sequence. In this case, we can find a direct bijection between the objects in Proposition 3.5 as described below.

We start by labeling the vertices of the (innermost) level $`n`$ of $`P_4^n`$ clockwise by 2, 23, 12, 123. The level $`n1`$ is labeled starting from the vertex immediately to the left of vertex labeled 2 as follows: $`4,45,[3]4,[3]45`$ (the meaning of brackets will be discussed below). More generally, for $`i<n`$, given a level $`ni+1`$ labeled clockwise with $`2i,2i(2i+1),[2i1]2i,[2i1]2i(2i+1)`$, we label the (next outer) level $`ni`$ clockwise from the inside out with $`2i+2,(2i+2)(2i+3),[2i+1](2i+2),[2i+1](2i+2)(2i+3)`$ starting from the vertex immediately to the left of vertex labeled $`2i`$. (See Figure 5 for the case $`n=3`$.) Each independent set has at most one vertex on each level. Now, given any independent set in $`P_4^n`$, we can write the labels of its nodes in increasing order and delete any integer $`[2i1]`$ in brackets if the sequence also contains $`2i2`$ or $`2i1`$ without brackets. Erasing all brackets now, if any, we obtain a sequence from Proposition 3.5.

###### Example 3.6.

The independent sets $`\{45,[5]67\},\{4,[5]67\},\{12,[5]6\},\{[3]4,67\}`$ correspond to the sequences $`4567,467,1256,3467`$, respectively.

For convenience, we will write down the set of nonadjacent labels at level $`ni`$ for each label at level $`ni+1`$.

(4)
$$\begin{array}{cc}\hfill \text{(no label) }ϵ& \left(2i+2\right),\left(2i+2\right)\left(2i+3\right),\left[2i+1\right]\left(2i+2\right),\left[2i+1\right]\left(2i+2\right)\left(2i+3\right)\hfill \\ \hfill 2i& \left[2i+1\right]\left(2i+2\right),\left[2i+1\right]\left(2i+2\right)\left(2i+3\right)\hfill \\ \hfill 2i\left(2i+1\right)& \left(2i+2\right),\left[2i+1\right]\left(2i+2\right)\left(2i+3\right)\hfill \\ \hfill \left[2i1\right]2i& \left(2i+2\right),\left(2i+2\right)\left(2i+3\right)\hfill \\ \hfill \left[2i1\right]2i\left(2i+1\right)& \left(2i+2\right)\left(2i+3\right),\left[2i+1\right]\left(2i+2\right)\hfill \end{array}$$

Now it is not difficult to construct an independent set given a sequence of Proposition 3.5. We partition the sequence of integers from 1 to $`2n+1`$ as follows:

$$123|\mathrm{\hspace{0.17em}45}|\mathrm{\hspace{0.17em}67}|\mathrm{}|\mathrm{\hspace{0.17em}2}n(2n+1),$$

then choose the vertices of the independent set in the order of increasing labels using the rules (4). Notice that the label of the vertex at level $`ni+1`$ must contain $`2i`$.

###### Example 3.7.

$$\begin{array}{cc}\hfill 4567& 123|\underset{¯}{45}|\underset{¯}{67}(ϵ,45,[5]67)\{45,[5]67\}\hfill \\ \hfill 467& 123|\underset{¯}{4}5|\underset{¯}{67}(ϵ,4,[5]67)\{45,[5]67\}\hfill \\ \hfill 1256& \underset{¯}{12}3|4\underset{¯}{5}|\underset{¯}{6}7(12,ϵ,[5]6)\{12,[5]6\}\hfill \end{array}$$

It can be shown that the two maps described above are inverses of each other based on the recursive structure of sequences under consideration.


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# Searching for sub-millisecond pulsars from highly polarized radio sources

## 1. Introduction

The NRAO VLA Sky Survey (NVSS; Condon et al. 1998) contains 1.8 million radio sources with flux density at 1400 MHz greater than 2.5 mJy with linear polarization information. Many pulsars are known to have high linear polarization on average. Han & Tain (1999) identified 97 known pulsars from the NVSS source catalogue and they found that, on average, the linear polarization percentage of pulsars is much higher than that of other classes of objects such as quasars and BL-Lac objects. Therefore, high linear polarization can be used as a criterion for selecting pulsars with all kinds of possible periods, including submillisecond pulsars for which normal untargetted surveys are presently impossible. Submillisecond pulsars could be strange-quark stars, which have been explored theoretically (e.g. Madsen 1998), but with no detection observationally yet (e.g. Edwards et al. 2001).

We have selected 253 unresolved sources with high linear polarization (linear polarization percentage $`L/S>25\%`$; the uncertainty of $`L/S<10\%`$) from the NVSS catalog and searched for pulsed emission. These sources have a flux density which is generally larger than $``$4 mJy. To have been missed by previous pulsar surveys, they must have some combination of the following properties: short (millisecond) pulse period, short orbital period, and/or high dispersion. Most of these sources are at relatively high Galactic latitudes and so the first two properties are likely to be the most important if they are pulsars.

We searched for short-period pulsed emission from these selected sources by using Parkes 64m telescope in October 8-19 1999 and December 9-11 1999. The central beam of the 20-cm multibeam system was used, having a system temperature of 21 K and a gain of 0.735 K Jy-1. The two orthogonally polarized signals were amplified and fed to a filter-bank system, which was used in one of two configurations: 512 channels of 0.5 MHz bandwidth centered at 1261.75 MHz, or 256 channels of 0.25 MHz bandwidth centered on 1293 MHz. After detection of each channel, the two polarization signals were added together and 1-bit digitized. Data were sampled every 0.08 ms for a total observation time of 340 seconds for the 512x0.5 MHz system, and 0.05 ms sample time for 450 seconds for the 256x0.25 MHz system.

Data-processing consisted of sub-band dedispersion, fine-dedispersion and periodicity search. As shown in Fig.1, our search is very sensitive to submillisecond pulsars up to a dispersion measure of 80 pc cm-3, much better than the Parkes 70 cm pulsar survey (Manchester et al. 1996) which had a sensitivity of 3 mJy at 400 MHz for long period pulsars, corresponding to about 0.5 mJy at 1400 MHz for a typical pulsar spectral index of $`1.5`$. That survey was insensitive to pulsars with period less than 2 ms for typical DMs.

No new pulsars were detected though test observations of known pulsars gave results as expected. This null result implies that either radio pulsars do not have pulse periods of less than 1 ms, or submillisecond pulsars are not highly polarized radio sources.

## References

Condon J.J., Cotton W.D., Greisen E.W., et al., 1998, AJ 115, 776

Edwards R. T., van Straten W., Bailes M., 2001, ApJ 560, 365

Han J.L., Tian W.W., 1999, A&AS 136, 571

Madsen J., 1998, Phy. Rev. Let. 81, 3311

Manchester R.N., Lyne A.G., D’Amico N., et al. 1996, MNRAS 279, 1235


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# Broken Symmetry in Density-Functional Theory: Analysis and Cure

## Abstract

We present a detailed analysis of the broken-symmetry mean-field solutions using a four-electron rectangular quantum dot as a model system. Comparisons of the density-functional theory predictions with the exact ones show that the symmetry breaking results from the single-configuration wave function used in the mean-field approach. As a general cure we present a scheme that systematically incorporates several configurations into the density-functional theory and restores the symmetry. This cure is easily applicable to any density-functional approach.

The nanoscale semiconductor systems are technically very promising for future components of microelectronic devices. From the theoretical point of view, quantum dot (QD) systems are valuable source of novel quantum effects. Many of these result from the fact that the electron-electron interaction and external magnetic field have greatly enhanced effects compared to atoms and molecules. This raises new challenges for the theoretical methods, and the validity of approximations in, e.g., mean-field approaches can be questioned. For this reason, QD systems serve as perfect test cases to develop the electronic structure methods, with the results still applicable to great variety of physical problems where mean-field approaches have been used.

In earlier studies, Hartree-Fock (HF) and especially density-functional theory (DFT) methods have shown to produce accurate results for various QD systems, even with small $`N`$. However, in the context of solutions with a broken spin symmetry, the validity of the mean-field approaches has been actively discussed in the literature RM . The spin-density wave (SDW) formation in QDs has been compared to similar phenomena found in isotropic metals overhauser , organic linear-chain compounds gruner , atomic nuclei frauendorf , and small fermion clusters hakkinen . According to the Jahn-Teller theorem, any non-linear molecular system in a degenerate electronic state becomes more stable by removing the degeneracy and thus lowering the symmetry and the total energy. A crucial difference between molecular and QD systems is, however, that as the nuclei in molecules are free to move and relax, the QD potential is external and fixed as it results from, e.g., metallic gates. Thus to lower the symmetry in QD, the spin densities must “relax” in an anti-ferromagnetic fashion to a SDW solution. This is claimed to reveal the electron correlations inherent in the true ground state RM .

In this Letter, we analyze symmetry breaking in a two-dimensional rectangular QD Esa  using both DFT and exact diagonalization (ED). We concentrate on the four-electron case, as it is the first particle number showing the general features of electronic structure seen also for larger particle numbers, such as the transitions between the two spin states $`S=0`$ and 1 and the SDW solution predicted by DFT. We find that SDW clearly reflects the limitations of basic DFT to describe systems that have more than one major configuration in the ground-state wave function. There is a continuous interest for developing DFT methods for this kind of systems. The main difficulty for DFT is the fact that these systems have ensemble-$`v`$-representable densities (E-VR) in contrast to the more common pure-state $`v`$-representable densities (P-VR) UK . As an interesting feature we see a continuous transition from an E-VR to a P-VR density as we deform our QD. Finally, we present a simple modification of DFT that is able to describe the multi-configurational nature of the ground states.

The generally used model Hamiltonian of an $`N`$-electron QD system can be written as

$$H=\underset{i=1}{\overset{N}{}}\left[\frac{\mathrm{}^2}{2m^{}}_i^2+V_{\mathrm{ext}}(𝐫_i)\right]+\underset{i<j}{\overset{N}{}}\frac{e^2}{ϵ|𝐫_i𝐫_j|},$$
(1)
where we have used the effective-mass approximation to describe electrons moving in the $`xy`$ plane, surrounded by background material of GaAs with the effective electron mass $`m^{}=0.067m_e`$ and dielectric constant $`ϵ=12.4`$. We use scaled atomic units, and energies are thus given in $`\mathrm{Ha}^{}11.86`$ meV and lengths in $`a_B^{}9.79`$ nm. The external confinement in the $`xy`$ plane is described by an infinite hard-wall potential,

$$V_{\mathrm{ext}}(x,y)=\{\begin{array}{cc}0,\hfill & 0x\beta L,\mathrm{\hspace{0.17em}0}yL\hfill \\ \mathrm{},\hfill & \text{ elsewhere}.\hfill \end{array}$$
(2)
The deformation parameter $`\beta `$ defines the ratio between the side lengths of the rectangle. The area of the dot is fixed to be $`\pi ^2`$. The single-particle eigenstates are sine-functions in both directions, labeled with two quantum numbers $`(n_x,n_y)`$, and energies $`E_{n_x,n_y}=\left(n_x^2/\beta +\beta n_y^2\right)/2`$. Fig. 1 shows the three lowest eigenstates and the most important $`S_z=0`$ configurations of the four-electron QD.

We solve the electronic structure of QD using ED and DFT. In ED, the many-particle wave function is constructed as an expansion of the non-interacting eigenstates. The results approach the exact ones as more terms are added to the expansion. We use a basis of up to around fifteen thousand configurations. The interaction matrix elements are calculated numerically using Gaussian integration. In ED, all many-body quantum effects are taken into account in an exact fashion. In DFT, these are incorporated in a mean-field fashion as an effective potential. In the DFT method used, we allow different spin densities for up and down electrons. This is necessary for $`S0`$, and needed also for $`S=0`$ in order to find broken-symmetry solutions. More details of the DFT method and the numerical implementation can be found from Ref. Esa  and references therein.

In Fig. 2 we present the DFT and ED energies of the rectangular quantum dot as a function of the deformation parameter $`\beta `$.

For $`\beta `$ close to unity, the $`S=1`$ state is lower in energy than the $`S=0`$ state, in accordance with Hund’s rule. In the case of the $`S=1`$ state, the DFT energies compare quite well with those obtained by ED: the deviation between the two remains nearly constant for all values of $`\beta `$. Such a behavior is not seen in the $`S=0`$ results, for which we show two DFT energies: one with retained symmetry and another with a broken one. The broken-symmetry solution has a non-zero total spin density, corresponding to a SDW solution, see Fig 3.

The DFT calculation with the spin symmetry does not converge for smaller $`\beta `$ than those shown. This is due to the degeneracy in the system. Convergence can be achieved by use of fractional occupations. Comparing the DFT $`S=0`$ energies to the exact ones, one can see that, unlike for the $`S=1`$ case, the error in DFT is not constant. The energy of the symmetry-restricted state grows linearly towards $`\beta 1`$ where the ED energy saturates. On the other hand, the SDW state has an energy that overcompensates the error in the symmetry-restricted energy. The energy of the SDW state is closer to the exact value than the energy of the proper-symmetry state. One should note that the errors in DFT energies nearly cancel at the ground-state transition point, and the DFT prediction for it is very accurate.

It is claimed that the SDW spin densities reflect the internal structure of the system RM . To analyze this claim, we have plotted the SDW spin density of the DFT and ED conditional densities in Fig. 3. The conditional density is defined to be the electron density of the remaining three electrons as the coordinates of one of the electrons are fixed. In addition, we plot the ED spin density for the sum of the $`S=0`$ and $`S=1`$ states. One can see from the densities that there is a clear anti-ferromagnetic order in the system. Densities for parallel spins are localized in the opposite corners. Apart from this fact, the similarity of the conditional densities to the SDW density of Fig. 3 is marginal. However, the similarity of the SDW density to the unphysical mixture of the two spin states is very clear. The only difference is that the DFT density is slightly more localized. One should note that this similarity of the SDW solution to a mixture of two different spin states is pointed out by Hirose and Wingreen using ED in restricted basis HW .

To understand the electronic structure of the system and to analyze the problem associated with the SDW solution, it is enough to consider only the most important configurations in the ED solution, presented in Fig 1. The $`S=1`$ state, the ground state for small $`\beta `$, has three different $`S_z`$ states which are degenerate in energy. The one with $`S_z=0`$ consists of configurations $`C_3`$ and $`C_4`$ with equal weights for all values of $`\beta `$. The $`S=0`$ state is the ground state for large $`\beta `$, and it consists of the configurations $`C_1`$ and $`C_2`$. For $`\beta =1`$, these have equal weights, but for larger $`\beta `$, $`C_2`$ moves higher in energy and has a smaller weight in the exact wave function. For $`\beta 1.2`$, $`C_1`$ is clearly the dominating configuration. One should note that at this value of $`\beta `$, proper symmetry is restored in the DFT solution. The most natural reason for the occurrence of the SDW solution is that the basic DFT is unable to take into account more than one important configuration for the construction of the Kohn-Sham orbitals and the resulting densities. In terms of the configurations, the DFT spin densities at $`\beta =1.2`$ correspond to $`C_1`$. For smaller $`\beta `$, however, the SDW spin densities can only be obtained by a linear combination of all four configurations. For $`\beta =1`$, this linear combination is equal to the unphysical mixture of two different spin states used for Fig. 3 (b) above.

It is possible to analyze the broken-symmetry solution more generally by considering a mean-field-type single-configuration wave function for two up and two down-spin electrons, occupying the orbitals $`\psi _0`$ and $`\mathrm{sin}(\theta _\sigma )\psi _1+\mathrm{cos}(\theta _\sigma )\psi _2`$, where $`\theta _\sigma `$ contains the variational freedom for a spin type $`\sigma `$. Expanding this wave function results in four configurations similar to $`\{C_i\}_{i=1}^4`$ above. Assuming a further similarity to the QD case for $`\beta =1`$, one can write a Hamiltonian matrix of the four configurations as

$$H=\left(\begin{array}{cccc}E_1& \delta & 0& 0\\ \delta & E_1& 0& 0\\ 0& 0& E_0& \delta \\ 0& 0& \delta & E_0\end{array}\right),$$
(3)
where the configurations couple via the off-diagonal matrix element $`\delta `$ (taken to be real). The four exact energies are $`E_0\pm \delta `$ and $`E_1\pm \delta `$. One can set without loss of generality $`E_0=0`$ and $`E_1=1`$. The single-configuration energies have an interesting dependence on $`\delta `$, shown in Fig. 4. We present the energy as a function of the two variational angles $`\theta `$ for cases $`\delta =0.2`$ and 0.8. For small $`\delta `$, the second orbital for the minimum-energy solution is $`\psi _1`$ for one spin type and $`\psi _2`$ for the other. For the case $`\delta =0.8`$, the minima are found with orbitals $`\psi _1+\psi _2`$ and $`\psi _1\psi _2`$. The resulting total wave function of this case can easily be found to be a sum of the two exact wave functions (with energies $`\delta `$ and $`1\delta `$), and the energy of the mean-field state is equal to the average of the two exact energies. Furthermore, if one assumes that $`\psi _1`$ has a node on the $`x`$ axis and $`\psi _2`$ on the $`y`$ one, one can find densities similar to the SDW solution above. Now for the QD, the value of $`\delta `$ is close to 0.8, and one can understand the occurrence of the SDW solution from this more general argument.

The SDW solutions of Ref. Esa  for larger particle numbers can equally well be understood based on the four-electron case and the general argument presented above. In the cases where two Kohn-Sham orbitals are degenerate, we have a $`S=1`$ ground state. When the aspect ratio $`\beta `$ is changed, the energies split and one always finds a broken-symmetry SDW solution. The similarity to the four-electron case follows from the fact that in all these cases, there are two spatial orbitals of both spin type occupied by two electrons. The spin density in the SDW structure can then be directly found from the two nearly degenerate states as $`(\psi _1+\psi _2)^2(\psi _1\psi _2)^2`$. For example, the densities in Fig. 6 of Ref. Esa  are accurately reproduced by this formula using for the degenerate states the non-interacting ones with quantum numbers $`(1,3)`$ and $`(3,2)`$ for the left panel, or $`(3,2)`$ and $`(4,1)`$ for the right panel.

Based on the results presented above, it is clear that standard DFT is not able to describe accurately E-VR systems. The method of Ullrich and Kohn UK  is one possible solution, but this method might have an underlying problem. Namely, even in the case of an open shell and degeneracy, there are systems that still are P-VR, simply because the configurations do not necessarily mix even if they are degenerate. One such example is the parabolic QD, where the angular momentum is a good quantum number and single-particle states can be chosen in such a fashion that only one major configuration is found. It is not straightforward to see how the method of Ref. UK  assort the open-shell cases that are E-VR from those that are still P-VR.

As a possible solution we propose a scheme where first a standard DFT calculation is performed for the system (without symmetry breaking but with fractional occupations for the degenerate levels). For our case with $`\beta =1`$, the occupations of the two highest orbitals $`\psi _1`$ and $`\psi _2`$ are $`1/2`$ and the DFT energy is 13.26. One can construct two $`S=S_z=0`$ DFT configurations that have density equal to the DFT one by defining new orbitals $`\varphi _\pm =\psi _1\pm i\psi _2`$. Now the configurations involve the core DFT orbital $`\psi _0`$, and either $`\varphi _+`$ or $`\varphi _{}`$. The occupied orbitals in configurations are the same for both spin electrons, similarly to $`C_1`$ and $`C_2`$ in Fig. 1. The coupling of these two DFT configurations can be approximated by

$$\delta =\varphi _+^{}(𝐫_1)\varphi _+^{}(𝐫_2)\frac{1}{r_{12}}\varphi _{}(𝐫_1)\varphi _{}(𝐫_2)𝑑𝐫_1𝑑𝐫_2.$$
(4)
Now the DFT energy gives the diagonal Hamiltonian matrix elements and $`\delta `$ the off-diagonal ones. The two-ensemble DFT energy can be found by diagonalizing the Hamiltonian matrix. For our example, the absolute value of delta is found to be $``$ 0.14. Thus the mixing of the two configurations lowers the DFT energy to 13.12. This value is consistent with the ED one, if one takes into account the difference in DFT and ED energies for the non-degenerate cases. One should note that for a parabolic QD, the absolute value of $`\delta `$ is zero (resulting from the rotational symmetry), meaning that the configurations (with different angular momentum) do not mix. This also shows that our scheme correctly predicts the system to be P-VR. A similar behavior can be found for the multiplets in open-shell atoms.

The underlying idea of the scheme presented above is that DFT is able to efficiently describe correlation effects in a certain subspace of the full Hilbert space. This sub-space is related to a one DFT configuration. In P-VR cases this is sufficient for the accurate description of the system, but for an E-VR case, there are two or more subspaces relevant for the ground state, and DFT is unable to couple these. This coupling can be introduced, and one natural way is via $`\delta `$ of Eq. (4) above.

The generalization of the scheme for cases without an exact degeneracy of the DFT orbitals is straightforward. In addition, the approximation made for $`\delta `$ can be directly used for cases with larger particle numbers, too. This is because the states that are occupied in both configurations do not appear in the formula for $`\delta `$. We believe that the presented approach shows to be useful for many applications of DFT, especially for molecules, where the calculations of chemical reactions have observed similar problems of basic DFT H2H2 . More details and comparisons with other ensemble DFT approaches are left for forthcoming studies.

Concluding, have shown that the use of a single-configuration wave function in a mean-field theory can lead to an unphysical solution with a broken symmetry. In our case of a four-electron rectangular QD, the energy of the SDW solution is reasonable, but the spin densities have only a minor similarity with the exact total or conditional ones. We also present an analysis with a more general Hamiltonian matrix and we feel that our findings are relevant for a great variety of systems studied by the mean-field approaches, DFT in particular. As a cure, we propose a scheme for incorporating systematically several configurations into a mean-field approach. The method presented avoids the necessity of symmetry breaking, and has a built-in criteria to determine if several configurations are actually needed or not.

###### Acknowledgements.

We thank R. van Leeuwen, V. Sverdlov, and E. Thuneberg for discussions and acknowledge support by the Academy of Finland’s Centers of Excellence Program (2000-2005).


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# Dynamics of Glue-Balls in 𝑁=1 SYM Theory

## 1 Introduction

The low energy effective action of $`N=1`$ SYM theory is written in terms of a chiral effective field $`S=\phi +\theta \psi +\theta ^2F`$, which may be defined from the local source extension of the SYM action

$`S`$ $`{\displaystyle \frac{\delta }{\delta J}}W[J,\overline{J}],`$ $`e^{iW[J,\overline{J}]}`$ $`={\displaystyle 𝒟Ve^{i{\scriptscriptstyle d^4xd^2\theta (J+\tau _0)\mathrm{Tr}W^\alpha W_\alpha }+\mathrm{h}.\mathrm{c}.}}.`$ (1)
With appropriate normalization $`S`$ is equivalent to the anomaly multiplet $`\overline{D}^{\dot{\alpha }}V_{\alpha \dot{\alpha }}=D_\alpha S`$. $`J(x)`$ is the chiral source multiplet, with respect to which a Legendre transformation can be defined . The resulting effective action is formulated in terms of the gluino condensate $`\phi \mathrm{Tr}\lambda \lambda `$, the glue-ball operators $`F\mathrm{Tr}F_{\mu \nu }F^{\mu \nu }+i\mathrm{Tr}F_{\mu \nu }\stackrel{~}{F}^{\mu \nu }`$ and a spinor $`\psi (\sigma ^{\mu \nu }\lambda )_\alpha F_{\mu \nu }`$. An effective Lagrangian in terms of this effective field $`S`$ has the form

$$_{\text{eff}}=\text{d}^4\theta K(S,\overline{S})\left(\text{d}^2\theta S(\mathrm{log}\frac{S}{\mathrm{\Lambda }^3}1)+\mathrm{h}.\mathrm{c}.\right).$$
(2)
The correct anomaly structure is realized by the superpotential and thus $`K(S,\overline{S})`$ is invariant under all symmetries. In ref.  the explicit ansatz $`K=k(\overline{S}S)^{1/3}`$ had been made, which leads to chiral symmetry breaking due to $`S=\mathrm{\Lambda }^3`$, but supersymmetry is not broken as $`\phi `$ and $`\psi `$ acquire the same mass $`m=\mathrm{\Lambda }/k`$.

## 2 Glue-balls and constraint Kähler geometry

Though the spectrum found in ref.  does not include any glue-balls, such fields do appear in $`F`$. However, they drop out in the analysis of , as $`F`$ is treated as an auxiliary field. Indeed, the highest component of a chiral superfield is auxiliary in standard SUSY non-linear $`\sigma `$-models, i.e. there appear no derivatives acting onto this field and moreover its potential is not bounded from below, but from above. In case of the Veneziano-Yankielowicz Lagrangian the part depending on the auxiliary field reads

$$_{\text{aux}}=k(\overline{\phi }\phi )^{\frac{2}{3}}\overline{F}F+\left(\frac{1}{3}\phi ^{\frac{2}{3}}\overline{\phi }^{\frac{5}{3}}F\overline{\psi }\overline{\psi }F\mathrm{log}\frac{\phi }{\mathrm{\Lambda }^3}+\mathrm{h}.\mathrm{c}.\right),$$
(3)
and the supersymmetric spectrum is obtained, *if and only if* $`F`$ is eliminated by the algebraic equations of motion that follow from (3). This leads to the unsatisfactory result that glue-balls cannot be introduced in a straightforward way (cf. also ) which, in addition, contradicts available lattice-data .

However, in the special case of $`N=1`$ SYM the elimination of $`F`$ is not consistent: If $`F`$ is eliminated from (3), this implies that the theory must be ultra-local in the field $`F`$ *exactly*, i.e. even corrections to the effective Lagrangian which are not included in (2) are not allowed to change the non-dynamical character of $`F`$. If this field would be related to the fundamental auxiliary field, this restriction would be obvious. But in $`N=1`$ SYM the situation is different: $`S`$ is the effective field from a composite operator and $`F`$ is not at all related to the fundamental auxiliary field $`D`$. As a consequence, the restriction of ultra-locality on $`F`$ leads to an untenable constraint on the *physical* glue-ball operators (for details we refer to ).

As shown in ref. , the effective Lagrangian of  is not the most general expression compatible with all the symmetries, but the constant $`k`$ may be generalized to a function $`k(\frac{S^{1/3}}{\overline{D}^2\overline{S}^{1/2}},\frac{\overline{S}^{1/3}}{D^2S^{1/2}})`$. This non-holomorphic part automatically produces space-time derivatives onto the field $`F`$, which is most easily seen when $`K(S,\overline{S})`$ is rewritten in terms of two chiral fields :

$$K(S,\overline{S})K(\mathrm{\Psi }_0,\mathrm{\Psi }_1;\overline{\mathrm{\Psi }}_0,\overline{\mathrm{\Psi }}_1)$$
(4)
$`\mathrm{\Psi }_0`$ and $`\mathrm{\Psi }_1`$ are not independent, but they must obey the constraints

$`\mathrm{\Psi }_0`$ $`=S^{\frac{1}{3}}=\phi ^{\frac{1}{3}}+{\displaystyle \frac{1}{3}}\phi ^{\frac{2}{3}}\theta \psi +{\displaystyle \frac{1}{3}}\theta ^2(\phi ^{\frac{2}{3}}F+{\displaystyle \frac{1}{3}}\phi ^{\frac{5}{3}}\psi \psi ),`$ (5)
$`\mathrm{\Psi }_1`$ $`=\overline{D}^2\overline{\mathrm{\Psi }}_0={\displaystyle \frac{1}{3}}(\overline{\phi }^{\frac{2}{3}}\overline{F}+{\displaystyle \frac{1}{3}}\overline{\phi }^{\frac{5}{3}}\overline{\psi }\overline{\psi }){\displaystyle \frac{i}{3}}\theta \sigma ^\mu _\mu (\overline{\phi }^{\frac{2}{3}}\overline{\psi })\theta ^2\mathrm{}\overline{\phi }^{\frac{1}{3}}.`$ (6)
As $`F`$ appears as lowest component of $`\overline{\mathrm{\Psi }}_1`$, the Lagrangian includes a kinetic term for that field. In contrast to the situation in , this is not inconsistent as the potential in $`F`$ may include arbitrary powers in that field (instead of a quadratic term only) and can be chosen to be bounded from below (instead of above). This way the field $`F`$ is promoted to a usual physical field. It has been shown in  that there exist consistent models of this type. In  these ideas have been applied to $`N=1`$ SYM, leading to an effective action of that theory with dynamical glue-balls as part of the low-energy spectrum. Formally, the effective potential looks the same as in the case of Veneziano and Yankielowicz:

$$\begin{array}{cc}\hfill V_{\text{eff}}& =\stackrel{~}{g}_{\phi \overline{\phi }}F\overline{F}+\frac{1}{2}\stackrel{~}{g}_{\phi \overline{\phi },\overline{\phi }}F(\overline{\psi }\overline{\psi })+\frac{1}{2}\stackrel{~}{g}_{\phi \overline{\phi },\phi }\overline{F}(\psi \psi )\frac{1}{4}\stackrel{~}{g}_{\phi \overline{\phi },\phi \overline{\phi }}(\psi \psi )(\overline{\psi }\overline{\psi })\hfill \\ & +c\left(F\mathrm{log}\frac{\phi }{\mathrm{\Lambda }^3}+\overline{F}\mathrm{log}\frac{\overline{\phi }}{\overline{\mathrm{\Lambda }}^3}\frac{1}{2\phi }(\psi \psi )\frac{1}{2\overline{\phi }}(\overline{\psi }\overline{\psi })\right)\hfill \end{array}$$
(7)
However, in contrast to  the Kähler “metric”111This quantity is not equivalent to the true Kähler metric of the manifold spanned by $`\mathrm{\Psi }_0`$ and $`\mathrm{\Psi }_1`$, cf. . is a function of $`\phi `$ *and* $`F`$, $`\stackrel{~}{g}_{\phi \overline{\phi }}(\phi ,F;\overline{\phi },\overline{F})`$. From eq. (7) the consistent vacua can be derived, for explicit expressions we refer to . The most important properties of the Lagrangian (2) with (4) are:

The effective potential is minimized with respect to *all* fields $`\phi `$, $`\psi `$ and $`F`$. Consequently, the dominant contributions that stabilize the potential must stem from the Kähler part, not from the superpotential: The superpotential is a holomorphic function in its fields and therefore its scalar part must have unstable directions. In the present context there exists no mechanism to transform these instabilities into stable but non-holomorphic terms.

Though the model has the same superpotential as the Lagrangian of ref.  its spectrum is completely different: Chiral symmetry breaks by a vacuum expectation value (vev) of $`\phi \mathrm{\Lambda }^3`$, but this mechanism is more complicated than in . Any stable ground-state must have non-vanishing vev of $`F`$. But $`F`$ is the order parameter of supersymmetry breaking and thus this symmetry is broken as well222The author of ref.  concluded that this model cannot have a stable *supersymmetric* ground-state. This is in agreement with our results, as the model breaks down as $`F0`$.. $`\psi `$ is a massless spinor, the Goldstino.

The supersymmetry breaking scenario is of essentially non-perturbative nature333The importance of such a breaking mechanism has been pointed out in  already, but a concrete description was not yet found therein.: it is not compatible with perturbative non-renormalization theorems, as the value of $`V_{\text{eff}}`$ in its minimum and the vev of $`T_{}^{\mu }{}_{\mu }{}^{}`$ are no longer equivalent. In particular, the former can be negative, while the latter is positively semi-definite due to the underlying current-algebra relations. To our knowledge this is the first model, where this type of supersymmetry breaking has found a concrete description (cf.  for details).

Any ground state with $`\stackrel{~}{g}_{\phi \overline{\phi }}0`$ can be equipped with stable dynamics for $`p^2<|\mathrm{\Lambda }|^2`$. In the construction of concrete kinetic terms it is important to realize that (4) may include expressions with explicit space-time derivatives. Again this is possible as $`F`$ is not interpreted as an auxiliary field.

In summary, the Lagrangian of ref.  is the most general one, which can be formulated in terms of the effective field $`S`$. Consistent ground-states can be found together with broken supersymmetry only. It would be interesting to compare these results with a *different* action, which has supersymmetric ground-states. But the ”pièce de résistance” for such an action is the fact, that it cannot start from the effective field $`S`$.

### Acknowledgements

The author would like to thank U. Ellwanger, J.-P. Derendinger, E. Kraus, P. Minkowski, Ch. Rupp and E. Scheidegger for interesting discussions. This work has been supported by the Austrian Science Foundation (FWF) project P-16030-N08.


text:
# On the nonexistence of conformally flat slices in the Kerr and other stationary spacetimes

## Abstract

It is proved that a stationary solutions to the vacuum Einstein field equations with non-vanishing angular momentum have no Cauchy slice that is maximal, conformally flat, and non-boosted. The proof is based on results coming from a certain type of asymptotic expansions near null and spatial infinity —which also show that the developments of Bowen-York type of data cannot have a development admitting a smooth null infinity—, and from the fact that stationary solutions do admit a smooth null infinity.

*Introduction.-* One of the basic tenets of what R. Penrose has called the establishment’s view concerning the Physics of black holes, is that an asymptotically flat spacetime containing black holes will eventually settle down to a Kerr black hole. This point of view has been supported by the *no hair* theorems of Israel, Carter and Robinson Isr67 ; Car71 ; Rob75b , and confirmed by numerical simulations and perturbative analysis of the evolution of black hole initial data with non-vanishing angular momentum —see e.g. BraSei95 ; BakBruCamLouTak01 ; DaiLouTak02 ; GleNicPriPul98 . A great majority of these numerical simulations have made use of initial data sets which are conformally flat, and have a second fundamental form given by the so-called Bowen-York Ansatz BowYor80 . As noted by several authors, the Bowen-York data do not reduce to Kerr data under any choice of the parameters. Now, there are several reasons why it would be of interest to have slices of the Kerr spacetime that are conformally flat. For example, one would like to analyse the evolution of a pair of initially very close spinning black holes as a perturbation of the Kerr solution, and not as a perturbation of Schwarzschild as it has been done so far —see e.g. GleNicPriPul98 .

Garat & Price GarPri00  have shown that the Kerr spacetime has no axially symmetric, conformally flat slices that smoothly approach to the standard Schwarzschildean slices of constant time as the angular momentum goes to zero. Their argument was based on a perturbative analysis of the Cotton-Bach tensor —which locally characterises conformally flat 3-dimensional hypersurfaces— of the prospective slices up to second order in the angular momentum. Nevertheless, as noted by the authors themselves, there is, in principle, room for the existence more exotic slices than the ones covered by their Ansatz.

This letter presents a theorem that precludes the existence of maximal, non-boosted, conformally flat slices not only in the Kerr spacetime, but also in any stationary spacetime with non-vanishing angular momentum. The proof is based in general conformal properties of the stationary solutions to the Einstein field equations, and makes use not only of the information contained in the slices, but also of features concerning their developments.

Let $`\stackrel{~}{S}`$ be a Cauchy hypersurface, and let $`\stackrel{~}{h}_{\alpha \beta }`$, $`\stackrel{~}{\chi }_{\alpha \beta }`$ be, respectively, its first and second fundamental forms. We shall write $`\stackrel{~}{\chi }=\stackrel{~}{\chi }_\gamma ^\gamma `$. The hypersurface $`\stackrel{~}{S}`$ could have several asymptotically flat regions. We concentrate our attention on one of them. Let $`\{y^\alpha \}`$ be coordinates in this asymptotic region such that,

$$\stackrel{~}{h}_{\alpha \beta }=\left(1+\frac{2m}{|y|}\right)\delta _{\alpha \beta }+𝒪\left(\frac{1}{|y|^2}\right),\stackrel{~}{\chi }_{\alpha \beta }=𝒪\left(\frac{1}{|y|^2}\right),$$

as $`|y|\mathrm{}`$, where $`m`$ is the ADM mass associated with this end. The (ADM) linear and angular momentum of the asymptotically flat region are defined by,

$`P_\alpha ={\displaystyle \frac{1}{8\pi }}\underset{r\mathrm{}}{lim}{\displaystyle _{S_r}}(\stackrel{~}{\chi }_{\alpha \beta }\stackrel{~}{\chi }\stackrel{~}{h}_{\alpha \beta })n^\beta 𝑑S_r,`$
$`J_\alpha ={\displaystyle \frac{1}{8\pi }}\underset{r\mathrm{}}{lim}{\displaystyle _{S_r}}ϵ_{\alpha \beta \gamma }y^\beta (\stackrel{~}{\chi }^{\gamma \rho }\stackrel{~}{\chi }\stackrel{~}{h}^{\gamma \rho })n_\rho 𝑑S_r,`$
respectively, where $`S_r`$ denotes the 2-sphere $`\{|y|=r\}`$, and $`n^\alpha `$ is the outward unit normal.

Our main theorem is the following:

###### Theorem.

Let $`(\stackrel{~}{S},\stackrel{~}{h}_{\alpha \beta },\stackrel{~}{\chi }_{\alpha \beta })`$ be a maximal, asymptotically Euclidean, conformally flat initial data set for the Einstein vacuum field equations. Assume, also, that the initial data set contains no linear momentum, and that the conformally rescaled second fundamental form admits an expansion of the form given by equation (2). If, furthermore, the spacetime is stationary, then the angular momentum of the data must vanish.

Thus, stationary spacetimes —whence the Kerr solution— admit no maximal, conformally flat, non-boosted slices with a second fundamental form of the sort prescribed by equation (2) if their ADM angular momentum is non-zero. If the angular momentum were non-zero, then the Schwarzschild solution would render a counterexample to our claim.

Some remarks concerning the hypothesis of the theorem come into place. Firstly, the assumption of maximal slices —i.e. such that $`\stackrel{~}{\chi }`$— is natural, for very little is know about how to solve the Einstein constraint equations for non-maximal data. Secondly, the slices have been assumed to be non-boosted, for the presence of linear momentum leads to non-smooth solutions of the Hamiltonian constraint —see DaiFri01 . Existence of solutions to the Hamiltonian constraint for a conformal second fundamental form of the sort given by equation (2) —under some extra conditions— has been discussed in the aforementioned reference.

*Sketch of the proof:* the proof of the theorem relies heavily on some conformal properties of stationary spacetimes and, in particular, on features of a certain kind of asymptotic expansions near spatial infinity that can be obtained from a representation of spatial infinity introduced by Friedrich Fri98a . For asymptotically Euclidean, conformally flat data with a second fundamental form like the one given in equation (2) —no stationarity being assumed here— , the asymptotic expansions contain certain logarithmic terms which preclude the existence of a smooth null infinity. These logarithmic terms vanish, if and only if, the angular momentum of the data vanishes. It can be seen that this implies that the development of such data admits no smooth null infinity. On the other hand, it is known that the developments of stationary data do admit a smooth null infinity. Consequently, if the initial data is to be stationary, its angular momentum must vanish. $`\mathrm{}`$

There are, nevertheless, a number of technical issues that one has to analyse. These will be discussed in the sequel. In particular, there is the fact that the gauge in which the expansions near spatial and null infinity are obtained is different to that in which the existence of a smooth null infinity for stationary spacetimes has been proved. It could well be the case that the non-smoothness in the expansions is a product of a bad choice of gauge. As it turns out, this is not the case, and the way one can do it is by comparing them in a gauge well adapted to null infinity —the NP gauge.

*Expansions near spatial and null infinity.-* Penrose Pen63  has introduced the seminal idea of studying the spacetimes describing isolated systems by means of a manifold with boundary –*the unphysical spacetime*$`(M,g_{\mu \nu })`$. This unphysical spacetime is obtained from the original one, $`(\stackrel{~}{M},\stackrel{~}{g}_{\mu \nu })`$ by means of a conformal rescaling $`g_{\mu \nu }=\mathrm{\Omega }^2\stackrel{~}{g}_{\mu \nu }`$. The boundary of the unphysical spacetime consists, at least, of the so-called null infinity $`I=I^+I^{}`$, and of a point, $`i^0`$, called spatial infinity. The spacetime $`(\stackrel{~}{M},\stackrel{~}{g}_{\mu \nu })`$ will be said to admit a smooth null infinity, if a conformal compactification can be introduced in such a way that the resulting null infinity is a smooth submanifold.

Friedrich has introduced a certain representation of spatial infinity as a cylinder —the cylinder at spatial infinity, $`I`$ Fri98a . A major feature of this representation of null infinity, is that the location and structure of the conformal boundary of the spacetime is known a priori, if the Einstein constraint equations have been solved. We shall refer to the gauge leading to this representation of spatial infinity as to the F-gauge. The equations governing the unphysical spacetime —the Conformal Einstein field equations— reduce, upon evaluation on $`I`$, to an interior system of transport equations. This system of transport equations enables us to obtain a certain type of asymptotic expansions, which in turn allow to connect properties of the spacetime at null infinity, with properties of the initial data. In particular, the components $`\varphi _j`$, $`j=0,\mathrm{},4`$ of the *rescaled Weyl spinor* $`\varphi _{abcd}`$ can be expanded as,

$$\varphi _j\underset{p0}{}\frac{1}{p!}\varphi _j^{(p)}(\tau ,\theta ,\phi )\rho ^p,$$
(1)
$`j=0,\mathrm{},4`$ where $`\tau `$ is a coordinate such that the locus of null infinity is given by $`\tau =\pm 1`$, and $`(\theta ,\phi )`$ are the standard polar coordinates. The cylinder at spatial infinity corresponds to the points for which $`\rho =0`$ and $`1\tau 1`$. The coefficients $`\varphi _j^{(p)}(\tau ,\theta ,\phi )`$ are determined by the interior system at $`I`$.

A set of scripts in the computer algebra system Maple V have been written in order to calculate the aforementioned expansions. They have been used to analyse the behaviour near spatial infinity of time symmetric solutions Val03b ; Val03c ; Val03e . Recently, these scripts have been modified to enable the discussion of conformally flat data with non-vanishing second fundamental form. The full results of these investigations will be presented elsewhere. Here, we shall only mention the results which are relevant for our discussion.

Consider a conformally flat, maximal initial data set. Under the so-called conformal Ansatz, the constraint equations reduce to:

$`\mathrm{\Delta }\vartheta ={\displaystyle \frac{1}{8}}\chi _{\alpha \beta }\chi ^{\alpha \beta }\vartheta ^7,`$
$`_\alpha (\vartheta ^4\chi ^{\alpha \beta })=0,`$
where $`\vartheta `$ is the conformal factor linking with the conformal metric, $`\delta _{\alpha \beta }=\vartheta ^4\stackrel{~}{h}_{\alpha \beta }`$. The conformally rescaled second fundamental form is given by $`\chi _{\alpha \beta }=\vartheta ^2\stackrel{~}{\chi }_{\alpha \beta }`$. The conformal factor $`\vartheta `$ can be used to compactify the initial hypersurface $`\stackrel{~}{S}`$. Let $`x^\alpha `$ denote normal coordinates such that $`\rho =|x|=0`$ at the point at infinity of the conformally rescaled initial hypersurface $`S`$.

In DaiFri01  Dain & Friedrich have discussed a the existence of initial data for the Einstein vacuum equations which near spatial infinity admit expansions of the form,

$$\stackrel{~}{h}_{\alpha \beta }\left(1+\frac{2m}{r}\right)\delta _{\alpha \beta }+\underset{k2}{}\frac{\stackrel{~}{h}_{\alpha \beta }^k}{r^k},\stackrel{~}{\chi }_{\alpha \beta }\underset{k2}{}\frac{\stackrel{~}{\chi }_{\alpha \beta }^k}{r^k},$$

where the coefficients involved in the expansions are smooth functions on $`S^2`$. This type of data is precisely of the sort one needs in order to use the CA scripts we have previously discussed. The sufficient conditions they found for the existence of such data imply that in the conformally flat case the conformally rescaled second fundamental form has to be of the form,

$$\chi _{\alpha \beta }=\chi _{\alpha \beta }^A+\chi _{\alpha \beta }^J+\chi _{\alpha \beta }^Q+𝒪(\rho ),$$
(2)
where

$`\chi _{\alpha \beta }^A={\displaystyle \frac{A}{\rho ^3}}\left(3n_\alpha n_\beta \delta _{\alpha \beta }\right),`$
$`\chi _{\alpha \beta }^J={\displaystyle \frac{3}{\rho ^3}}\left(n_\beta ϵ_{\gamma \alpha \rho }J^\rho n^\gamma +n_\alpha ϵ_{\rho \beta \gamma }J^\gamma n^\rho \right),`$
$`\chi _{\alpha \beta }^Q={\displaystyle \frac{3}{2\rho ^2}}\left(Q_\alpha n_\beta +Q_\beta n_\alpha (\delta _{\alpha \beta }n_\alpha n_\beta )Q^\gamma n_\gamma \right),`$
and $`n^\alpha =x^\alpha /\rho `$ is the radial unit normal near infinity —see also BeiOMu96 . The constants $`A`$, $`J^\alpha `$, $`Q^\alpha `$ are associated to the conformal Killing vectors on the Euclidean 3-dimensional space generating dilatations, rotations and “special conformal transformations” respectively.

For non-time symmetric data with a second fundamental form of the sort given in equation (2), the coefficients $`\varphi _j^{(0)}`$, $`\varphi _j^{(1)}`$, $`\varphi _j^{(2)}`$, and $`\varphi _j^{(3)}`$ in the expansion (1) are smooth functions of their arguments. However, for $`\varphi _j^{(4)}`$ one has that,

$$\varphi _j^{(4)}=f_j^{\mathrm{}}+f_j^{}\mathrm{ln}(1+\tau )+f_j^+\mathrm{ln}(1\tau ),$$

where $`f_j^{\mathrm{}}`$, $`f_j^\pm `$ are analytic functions of $`(\tau ,\theta ,\phi )`$. Furthermore, $`f_j^\pm `$ are polynomials in $`\tau `$ —with coefficients depending on $`(\theta ,\phi )`$— of the form,

$`f_0^+(\tau )=m|J|^4P_4(\tau )(1\tau )^2,`$ (3a)
$`f_1^+(\tau )=m|J|^3P_3(\tau )(1\tau )^3,`$ (3b)
$`f_2^+(\tau )=m|J|^2P_2(\tau )(1\tau )^4,`$ (3c)
$`f_3^+(\tau )=m|J|^2P_1(\tau )(1\tau )^5,`$ (3d)
$`f_4^+(\tau )=m|J|^2(1\tau )^6,`$ (3e)
and $`f_j^{}(\tau )=f_j^+(\tau )`$. We have also written $`|J|^2=J_\alpha J^\alpha `$. The functions $`P_1(\tau ),\mathrm{},P_4(\tau )`$ are polynomials in $`\tau `$ such that $`P_j(\pm 1)0`$. Thus, for conformally flat initial data, the presence of angular momentum is incompatible with a smooth null infinity. Note, that if the initial data is not conformally flat —as in the case of the standard Kerr data in Boyer-Lindquist coordinates— the situation is bound to be completely different. In this case the non-conformal parts of the metric would, in principle, leave room to cancel out the logarithmic singularities, although this must not be necessarily the case.

In order to compare the expansions implied by the aforementioned results, it is convenient to change our gauge from the F-gauge to the NP gauge. This is achieved by means of a conformal rescaling and a rotation of the tetrad in which the conformally rescaled Weyl spinor $`\varphi _{abcd}`$ is expressed. A full detailed list of what are the requirements for the NP gauge is given in FriKan00 . After a lengthy calculation, it can be shown that the expansions (1) together with (3a)-(3e) imply in the NP gauge the following expansions for the components of the Weyl spinor, $`\stackrel{~}{\mathrm{\Psi }}_{abcd}`$, of the physical (unrescaled) spacetime,

$`\stackrel{~}{\mathrm{\Psi }}_0k_0m|J|^2\mathrm{\Omega }^3\mathrm{ln}\mathrm{\Omega }+\stackrel{~}{\mathrm{\Psi }}_0^3\mathrm{\Omega }^3\mathrm{},`$ (4a)
$`\stackrel{~}{\mathrm{\Psi }}_1k_1m|J|^2\mathrm{\Omega }^3\mathrm{ln}\mathrm{\Omega }+\stackrel{~}{\mathrm{\Psi }}_1^3\mathrm{\Omega }^3\mathrm{},`$ (4b)
$`\stackrel{~}{\mathrm{\Psi }}_2k_2m|J|^2\mathrm{\Omega }^3\mathrm{ln}\mathrm{\Omega }+\stackrel{~}{\mathrm{\Psi }}_2^3\mathrm{\Omega }^3\mathrm{},`$ (4c)
$`\stackrel{~}{\mathrm{\Psi }}_3\stackrel{~}{\mathrm{\Psi }}_3^2\mathrm{\Omega }^2+\mathrm{},`$ (4d)
$`\stackrel{~}{\mathrm{\Psi }}_4\stackrel{~}{\mathrm{\Psi }}_4^1\mathrm{\Omega }+\mathrm{}.`$ (4e)
Here $`k_0`$, $`k_1`$ and $`k_2`$ denote non-zero constants, while the coefficients $`\stackrel{~}{\mathrm{\Psi }}_j^n`$ are functions of a retarded time $`u`$, which is defined up to a supertranslation, and the angular coordinates $`(\theta ,\phi )`$. Thus, note that if the angular momentum does not vanish, then null infinity not only is non-smooth, but the *Peeling behaviour* is not satisfied. It is noted that a somewhat similar result has been obtained by Klainerman and Nicolò —cfr. the remark of theorem 1.2 in KlaNic03 . It is also noted that the quantity $`m|J|^2`$ is essentially the constants associated to the leading logarithmic terms which have been discussed in ChrMacSin95 ; Val99a . The details of this, and a further connection to the so-called logarithmic Newman-Penrose constants will be also discussed elsewhere.

It is also noted that the above expansions imply that the development of Bowen-York type data with non-vanishing ADM angular momentum does not admit a smooth null infinity. This statement holds for the data constructed by Brandt & Brügmann BraBru97  as they are also conformally flat and use the same Ansatz for the second fundamental form as the considered by Bowen & York.

*The stationary solutions near null infinity.-* That the stationary solutions can be analytically extended through null infinity is known from the appendix in DamSch90  and the results of Dai01b . Now, the metric given in DamSch90  is not, in principle, in the NP gauge. Consider a null frame $`\{l^\mu ,n^\mu ,m^\mu ,\overline{m}^\mu \}`$ for that metric. A lengthy, but straightforward argument shows that the coordinate transformations and the Lorentz transformations required to take Damour & Schmidt’s metric to the NP gauge are analytic functions. Indeed, one has the following,

###### Lemma.

The stationary solutions written in the NP gauge have an analytic null infinity.

Moreover, it can be shown that —see e.g. NewPen68 ; PenRin86 —,

$`\stackrel{~}{\mathrm{\Psi }}_0=\stackrel{~}{\mathrm{\Psi }}_0^5\mathrm{\Omega }^5+𝒪(\mathrm{\Omega }^6),`$ (5a)
$`\stackrel{~}{\mathrm{\Psi }}_1=\stackrel{~}{\mathrm{\Psi }}_1^4\mathrm{\Omega }^4+𝒪(\mathrm{\Omega }^5),`$ (5b)
$`\stackrel{~}{\mathrm{\Psi }}_2=m\mathrm{\Omega }^3+𝒪(\mathrm{\Omega }^4),`$ (5c)
$`\stackrel{~}{\mathrm{\Psi }}_3=𝒪(\mathrm{\Omega }^3),`$ (5d)
$`\stackrel{~}{\mathrm{\Psi }}_4=𝒪(\mathrm{\Omega }^3),`$ (5e)
where the coefficients $`\stackrel{~}{\mathrm{\Psi }}_j^n`$ are $`u`$-independent.

Thus, one sees that if the further requirement of stationarity is imposed on the expansions (4a)-(4e) then necessarily $`|J|^2=0`$.

The only property of the stationary solutions that has been used is that they admit a smooth null infinity. This may look at first sight rather odd. However the results obtained in Val03e ; Val03c  and the ones discussed here suggest that admitting a smooth null infinity near spatial infinity may be an alternative characterisation of the stationary spacetimes. This, nevertheless remains to be proved.

###### Acknowledgements.

I would like to thank R. Beig, M. Mars, S. Dain and R. Lazkoz for very helpful discussions and suggestions. This project is funded by a Lise Meitner fellowship (M690-N09) of the Fonds zur Forderung der Wissenschaftlichen Forschung —FWF— Austria. The CA calculations described have been carried out in the computers of the Max Planck Institut für Gravitationsphysk, Albert Einstein Institut, Golm bei Potsdam, Germany.