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# Root games on Grassmannians

## 1 Prior work

In \[P\] we introduced the root game, a combinatorial game which can often determine whether or not a given Schubert structure constant is zero in the cohomology ring of a generalised flag manifold $`G/B`$. Our goal in this paper is to strengthen our earlier results in the case where the group $`G`$ is $`GL(n,)`$ and the Schubert classes are pulled back from a Grassmannian. We begin by recalling the root game for Schubert intersection numbers on the (ordinary) flag manifold $`Fl(n)`$. Apart from the root game, most of the relevant background material for this paper can be found in \[F1, FG\].

Let $`G=GL(n)`$. Let $`B`$ and $`B_{}`$ denote the Borel subgroups of upper and lower triangular matrices respectively, and $`T=BB_{}`$ the standard maximal torus, consisting of invertible diagonal matrices.

Recall that for each element of the symmetric group $`\pi S_n`$, there is a corresponding $`T`$-fixed point $`\pi B`$ on the flag manifold $`Fl(n)=G/B`$ (here we view $`\pi `$ as an element of $`GL(n)`$ via the standard representation of $`S_n`$), and an associated Schubert variety $`X_\pi =\overline{B_{}\pi B}`$. We denote its cohomology class in $`H^{}(Fl(n))`$ by $`[X_\pi ]`$.

Our convention will be to write all permutations in one line notation

$$\pi =\pi (1)\pi (2)\mathrm{}\pi (n).$$

If $`1S_n`$ denotes the identity element, and $`w_0=n\mathrm{}321S_n`$ is the long word, then the Schubert class $`[X_1]`$ is the identity element in $`H^{}(G/B)`$, and $`[X_{w_0}]H^{\text{top}}(G/B)`$ is the class of a point. In general $`[X_\pi ]`$ is a class of degree $`2\mathrm{}(\pi )`$, where $`\mathrm{}(\pi )`$ denotes the length of $`\pi `$.

For $`\pi _1,\mathrm{},\pi _mS_n`$, the Schubert intersection number

$$_{Fl(n)}[X_{\pi _1}]\mathrm{}[X_{\pi _m}]$$
(1)
is always a non-negative integer. The root game attempts to determine whether this number is strictly positive.

The game is played on a set of squares $`𝒮=\{S_{ij}|1i<jn\}`$. In our diagrams, we will arrange the squares $`S_{ij}`$ in an array, where $`i`$ is the row index and $`j`$ is the column index. In each square we allow tokens to appear. Each token has a label $`k\{1,\mathrm{},m\}`$, and no square may ever contain two tokens with the same label. A token labelled $`k`$ is called a $`k`$-token, and we write $`kS_{ij}`$ if a $`k`$-token appears in square $`S_{ij}`$.

A position in the game is specified by two pieces of data:

* The configuration of the tokens. Formally this is a map $`\tau `$ from $`𝒮`$ to subsets of $`\{1,\mathrm{},m\}`$, and our notation $`kS`$ is shorthand for $`k\tau (S)`$; however, in this paper we will wish to think of each token as a physical object which can be moved from square to square.
* A partition of the set of squares $`𝒮=R_1\mathrm{}R_r`$. Each $`R_i`$ is called a region.

The initial position of the game is as follows: there is a single region $`R_1=𝒮`$, and for $`i<j`$, a $`k`$-token appears in square $`S_{ij}`$ if and only if $`\pi _k(i)>\pi _k(j)`$.

From the initial position we move the tokens in the manner prescribed in the next paragraph. However, before each move we have the option of splitting regions into multiple regions. We define an ideal subset of the squares to be a set $`A𝒮`$ with the property that if $`S_{ij}A`$, $`i^{}i`$ and $`j^{}j`$ then $`S_{i^{}j^{}}A`$. Given an ideal subset of the squares we can break up a region $`R`$ into two regions: $`RA`$ and $`RA`$. We call this splitting $`R`$ along $`A`$, and we may repeat the process as many times as desired.

A move is specified by a region $`R`$, a token label $`k\{1,\mathrm{},m\}`$ and a pair $`(i,j)`$ with $`1i<jn`$. After choosing these data, we move tokens as follows:

* For every $`h`$ with $`j<hn`$, if $`S_{jh}`$ and $`S_{ih}`$ are both in $`R`$ and a $`k`$-token appears in $`S_{jh}`$ but not in $`S_{ih}`$, we move the $`k`$-token from $`S_{jh}`$ to $`S_{ih}`$;
* For every $`h`$ with $`1h<i`$, if $`S_{hi}`$ and $`S_{hj}`$ are both in $`R`$ and a $`k`$-token appears in $`S_{hi}`$ but not in $`S_{hj}`$, we move the $`k`$-token from $`S_{hi}`$ to $`S_{hj}`$.

More succinctly put, within the region $`R`$ we move $`k`$-tokens horizontally from column $`i`$ to column $`j`$ and vertically from row $`j`$ to row $`i`$, wherever possible. See Figure 1.

In the play of the game we may make any sequence of moves in any order. The game is won when there is exactly one token in each square.

###### Remark 1.1.

It turns out to be advantageous to split along an ideal subset $`A`$ if and only if the total number of tokens in all squares in $`A`$ equals $`|A|`$. We call the process of finding all such $`A`$ and splitting along them splitting maximally. Although in this paper we won’t take full advantage of this fact by splitting maximally before every move, we will never even consider the possibility of splitting along $`A`$ if this condition is not met.

The main result that we shall need is the following.

###### Theorem 1 (\[P\]).

If the game for $`\pi _1,\mathrm{},\pi _m`$ can be won, then

$$_{Fl(n)}[X_{\pi _1}]\mathrm{}[X_{\pi _m}]1.$$

In general we do not know if the converse of Theorem 1 is true. When $`m=3`$, the Schubert intersection numbers (1) are structure constants for the ring $`H^{}(Fl(n))`$, and in this case the converse has been confirmed for $`n7`$. It would certainly be remarkable if it were true in general.

We can also use the game to study the cohomology rings of partial flag manifolds by pulling back cohomology classes to the full flag manifold. In this paper, we investigate this in some detail in the case of the Grassmannian. Our main result is a version of the converse of Theorem 1 for Grassmannian Schubert calculus.

When $`m=3`$ and the classes come from a Grassmannian, the intersection numbers (1) are Littlewood-Richardson numbers. These numbers are also important in representation theory and in the theory of symmetric functions—they are the structure constants for the representation ring of $`GL(n)`$, and for the ring of symmetric functions in the Schur basis. As such, they are well studied, and there are a number of combinatorial rules, and geometric rules (e.g. \[C, V\]) known for computing these numbers.

###### Remark 1.2.

An interesting and pleasant feature of root games is that the rules naturally extend to any number Schubert classes in a way which is manifestly symmetrical in these inputs: it is immediately clear from the definitions that reordering the input permutations $`\pi _1,\mathrm{},\pi _m`$ does not affect whether or not the game can be won. This property is not generally shared by other combinatorial or geometric rules for Littlewood-Richardson numbers. Most of the known rules have no manifest symmetry. Knutson-Tao puzzles \[KTW\] are perhaps the most manifestly symmetrical Littlewood-Richardson rule, having a $`3`$-cyclic symmetry when $`m=3`$, but this is lost when one attempts to generalise beyond beyond triple intersections.

The manifest symmetry of the root game on $`Fl(n)`$, descends to the Grassmannian case. In Section 2 we partially break this symmetry, but as we explain in Section 4, the full symmetry is easily restored. To our knowledge, the only other manifestly symmetrical rule in Grassmannian Schubert calculus is the Horn recursion (see the survey article \[F2\]), which, like the root game, does not explicitly compute Littlewood-Richardson numbers and is only for determining which Schubert intersection numbers are strictly positive.

## 2 Associating a game to a Grassmannian Schubert calculus problem

###### Definition 2.1.

A $`01`$-string is a string $`\sigma =\sigma ^1\mathrm{}\sigma ^n`$ where each $`\sigma ^i\{0,1\}`$. A $`0^m1^l`$-string is a $`01`$-string $`\sigma =\sigma ^1\mathrm{}\sigma ^{m+l}`$, where exactly $`l`$ of the $`\sigma ^i`$ are equal to $`1`$.

Schubert varieties in the Grassmannian $`Gr_l(n)`$ are indexed by $`0^{nl}1^l`$-strings. Fix a base flag

$$\{0\}=V_0V_1\mathrm{}V_n=^n$$

in $`^n`$. The Schubert variety $`Y_\sigma Gr_l(n)`$ corresponding to $`\sigma `$ is

$$Y_\sigma =\{y^n|dimyV_i\sigma ^1+\mathrm{}+\sigma ^i\}.$$

We denote its cohomology class in $`H^{}(Gr_l(n))`$ by $`[Y_\sigma ]`$. According to these conventions $`[Y_{0\mathrm{}01\mathrm{}1}]`$ is the identity element in $`H^{}(Gr_l(n))`$ and $`[Y_{1\mathrm{}10\mathrm{}0}]H^{\text{top}}(Gr_l(n))`$ is the class of a point.

Given a list of $`s+2`$ $`0^{nl}1^l`$-strings $`\sigma _1,\mathrm{},\sigma _s,\mu ,\nu `$, we will wish to study the Grassmannian Schubert intersection numbers

$$_{Gr_l(n)}[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ].$$

We do so by investigating an equivalent problem on a full flag manifold. The most obvious way to do this is to simply consider the product of the classes $`\alpha ^{}([Y_{\sigma _i}])`$ etc., under the natural map $`\alpha :Fl(n)Gr_l(n)`$, as in Lemma 2.2 below. However, our purposes require that we do things in a somewhat less straightforward way.

Let $`N0`$ be an integer, and let $`\sigma `$ be a $`0^{nl}1^l`$-string. Let $`i_1<\mathrm{}<i_{nl}`$ denote the positions of the zeroes in $`\sigma `$, and $`j_1<\mathrm{}<j_l`$ denote the positions of the ones. We define three ways to associate a permutation to the $`01`$-string $`\sigma `$:

$`\pi (\sigma ,N)`$ $`=i_1\mathrm{}i_{nl}j_1\mathrm{}j_l(n+1)(n+2)\mathrm{}(n+N)`$
$`\pi ^{}(\sigma ,N)`$ $`=(i_1+N)\mathrm{}(i_{nl}+N)12\mathrm{}N(j_1+N)\mathrm{}(j_l+N)`$
$`\pi ^{\prime \prime }(\sigma ,N)`$ $`=i_{nl}\mathrm{}i_1(n+N)\mathrm{}(n+1)j_l\mathrm{}j_1.`$
From $`\sigma _1,\mathrm{},\sigma _s,\mu ,\nu `$ we produce a list of permutations, $`\pi _1,\mathrm{},\pi _{s+2}S_{n+N}`$:

$`\pi _1`$ $`=\pi (\sigma _1,N)`$
$`\mathrm{}`$
$`\pi _s`$ $`=\pi (\sigma _s,N)`$
$`\pi _{s+1}`$ $`=\pi ^{}(\mu ,N)`$
$`\pi _{s+2}`$ $`=\pi ^{\prime \prime }(\nu ,N).`$
###### Proposition 2.1.

$$_{Gr_l(n)}[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]=_{Fl(n+N)}[X_{\pi _1}]\mathrm{}[X_{\pi _{s+2}}].$$

The proof is based on the following standard pullback calculations, whose proofs we omit.

###### Lemma 2.2.

Let $`\alpha :Fl(n)Gr_l(n)`$ be the map which forgets all but the $`l`$-dimensional subspace of the flag. Then $`\alpha ^{}([Y_\sigma ])=[X_{\pi (\sigma ,0)}]`$.

For a $`01`$-string $`\sigma `$, let $`\sigma _+`$ denote the string $`\sigma `$ followed by $`N`$ ones, and let $`{}_{+}{}^{}\sigma `$ denote the string $`\sigma `$ preceded by $`N`$ ones.

###### Lemma 2.3.

Let $`\beta :Gr_l(n)Gr_{l+N}(n+N)`$ be the map $`VV\times ^N^n\times ^N`$. Then $`\beta ^{}([Y_{\sigma _+}])=[Y_\sigma ]`$. If $`\sigma ^{}`$ is not of the form $`\sigma _+`$ for some $`0^{nl}1^l`$-string $`\sigma `$ then $`\beta ^{}([Y_\sigma ^{}])=0`$.

If $``$ is a partial flag variety and $`hH^{}()`$ is a Schubert class, let $`h^{}`$ denote the opposite Schubert class, i.e. the unique Schubert class such that $`_{}hh^{}=1`$. For example $`[X_\pi ]^{}=[X_{w_0\pi }]`$, and $`[Y_\sigma ]^{}=[Y_{\sigma ^{\text{rev}}}]`$ where $`\sigma ^{\text{rev}}=\sigma ^n\mathrm{}\sigma ^1`$ is $`\sigma `$ reversed.

If $`h_1,\mathrm{},h_rH^{}()`$ are Schubert classes then the statement that

$$_{}h_1\mathrm{}h_r=c$$
(2)
is equivalent to the statement that

$$h_1\mathrm{}\widehat{h_i}\mathrm{}h_r=ch_i^{}+\mathrm{}$$
(3)
in the Schubert basis. We’ll call Equation (3) the $`h_i`$-special version of Equation (2).

###### Proof of Proposition 2.1.

Consider the equation in $`H^{}(Gr_{l+N}(n+N))`$:

$$_{Gr_{l+N}(n+N)}[Y_{\sigma _{1}^{}{}_{+}{}^{}}]\mathrm{}[Y_{\sigma _{s}^{}{}_{+}{}^{}}][Y_{{}_{+}{}^{}\mu }][Y_{\nu _+}]=c.$$
(4)
If we take the $`[Y_{{}_{+}{}^{}\mu }]`$-special version of Equation (4) and pull it back to $`Gr_l(n)`$, we get (using Lemma 2.3) the $`[Y_\mu ]`$-special version of

$$_{Gr_l(n)}[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]=c.$$

On the other hand, if we take the $`[Y_{\nu +}]`$-special version of Equation (4) and pull it back to $`Fl(n+N)`$, we get (using Lemma 2.2) the $`[X_{\pi _{s+2}}]`$-special version of

$$_{Fl(n+N)}[X_{\pi _1}]\mathrm{}[X_{\pi _{s+2}}]=c.$$

## 3 Non-vanishing for Grassmannians

Let $`\sigma _1,\mathrm{},\sigma _s,\mu ,\nu `$ be $`0^{nl}1^l`$-strings (if $`s=1`$, we’ll write $`\sigma `$ instead of $`\sigma _1`$). For any given $`N0`$, we associate permutations $`\pi _1,\mathrm{},\pi _{s+2}`$ as before. Our goal in this section is to prove the following theorem.

###### Theorem 2.

Take $`N`$ suitably large ($`Nl`$ will always suffice). The root game corresponding to $`\pi _1,\mathrm{},\pi _{s+2}`$ can be won if and only if

$$_{Gr_l(n)}[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]1.$$

Moreover, only moves involving tokens labelled $`1,\mathrm{},s`$ are required.

###### Remark 3.1.

The sufficient condition $`Nl`$ is not a sharp bound. This fact raises a number of interesting questions, which will be discussed in Sections 4 and 5. In the meantime the reader should not be alarmed by examples which use smaller values of $`N`$.

We shall first consider what happens in the case where $`s=1`$.

Recall the correspondence between $`01`$-strings and Young diagrams. Our Young diagrams will be in the French convention (the rows are left justified and increase in length as we move down). If $`\sigma `$ is a $`01`$-string $`\sigma `$, let $`r_i(\sigma )`$ denote the number of ones before the $`i^{\mathrm{th}}`$ zero. We associate to $`\sigma `$ the Young diagram $`\lambda (\sigma )`$ whose $`i^{\mathrm{th}}`$ row is $`r_i(\sigma )`$, where we are allowing the possibility that some rows may have length $`0`$.

If $`\lambda `$ is a Young diagram, let $`N+\lambda `$ denote the Young diagram obtained by adding $`N`$ squares to each row of $`\lambda `$ including those rows which contain $`0`$ squares.

The initial positions of the $`1`$-tokens are in the shape of the Young diagram $`\lambda _1=\lambda (\sigma )`$. The initial positions of the $`2`$-tokens are in the shape of a Young diagram $`\lambda _2=N+\lambda (\mu )`$. The squares that do not contain a $`3`$-token are also in the shape of a Young diagram, which we’ll denote $`\lambda _{\overline{3}}`$; viewed upside down, $`\lambda _{\overline{3}}`$ is the complement to $`\lambda (\nu )`$ inside an $`(nl)\times (l+N)`$ rectangle. The lower left corner of each of the Young diagrams $`\lambda _1`$, $`\lambda _2`$, $`\lambda _{\overline{3}}`$ is in the square $`S_{nl,nl+1}`$. See Figure 2 for an illustration of how these shapes are generated.

If $`\lambda _2\lambda _{\overline{3}}`$, then it is a basic fact that $`[Y_\mu ][Y_\nu ]=0H^{}(Gr_l(n))`$ (in fact this is a necessary and sufficient condition). Therefore, we may assume that no square contains both a $`2`$-token and a $`3`$-token. The squares which contain neither a $`2`$-token nor a $`3`$-token are empty squares—since $`N`$ is suitably large the $`1`$-tokens are all to the left of these squares— and are in the shape of a skew-diagram $`\lambda _{\overline{3}}/\lambda _2`$.

At the outset of the game, some immediate splitting can occur. We split in such a way that each square containing a $`3`$-token becomes a $`1`$-square region of its own. (For some choices of $`(\sigma ,\mu ,\nu )`$ it may be possible to split beyond this, but our argument is slightly simplified if we elect not to.) The remaining squares are those of $`\lambda _{\overline{3}}`$, which form what we call the big region. The big region is the only region which is unsolved; naturally, therefore, this will be the region in which all moves take place.

It is worth taking a moment to note how tokens move within the big region. A priori, a move $`(i,j)`$ will cause some $`k`$-tokens to move horizontally and others to move vertically. However, since the rows of the big region are indexed by $`\{1,\mathrm{},nl\}`$, and the columns are indexed by the disjoint set $`\{nl+1,\mathrm{},n+N\}`$, these cannot both happen. If $`i<jnl`$ then $`k`$-tokens will move vertically from row $`j`$ to row $`i`$. If $`nl+1i<j`$ then $`k`$-tokens move horizontally from column $`i`$ to column $`j`$. No tokens move if $`inl<j`$.

###### Definition 3.2 (Zelevinsky \[Z\]).

A picture between two (French) skew diagrams is a bijection between their boxes with the property that if box $`A`$ is weakly above and weakly right of box $`B`$ in one diagram, then the corresponding boxes $`A^{}`$ and $`B^{}`$ are in lexicographic order (i.e. $`S_{ij}`$ precedes $`S_{i^{}j^{}}`$ if $`i<i^{}`$ or $`i=i^{}`$ and $`j<j^{}`$) in the other diagram.

Now $`_{Gr_l(n)}[Y_\sigma ][Y_\mu ][Y_\nu ]`$ is given by the Littlewood-Richardson coefficient $`c_{\lambda _1\lambda _2}^{\lambda _{\overline{3}}}`$, which can be described as the number of pictures between $`\lambda _1`$ and $`\lambda _{\overline{3}}/\lambda _2`$ \[Z\]. This is a reformulation of the Littlewood-Richardson rule \[LR\], closely related to the Remmel-Whitney formula \[RW\]. In particular, if this number is non-zero, there exists such a picture. We pick one, and denote by $`f`$ the map it defines from the squares of $`\lambda _1`$ to the squares of $`\lambda _{\overline{3}}/\lambda _2`$. Note that all of these squares are in the big region of the game.

To show that the game can be won, we will give an algorithm— the Grassmannian root game algorithm (GRGA)—which uses $`f`$ to construct a sequence of moves that transports each $`1`$-token to a square of $`\lambda _{\overline{3}}/\lambda _2`$.

Essential to the GRGA is the following numbering scheme. At each point in the game we associate a number—called the readiness number—to each unplaced $`1`$-token (i.e. one which has not already reached its final destination) and each empty square of $`\lambda _{\overline{3}}/\lambda _2`$.

###### Definition 3.3.

At any point in the game, let $`t`$ be a $`1`$-token whose initial square was $`S\lambda _1`$ and whose current square is $`S_{ij}`$. Let $`S_{i^{}j^{}}=f(S)`$. Define the readiness number of both the token $`t`$ and the square $`S_{i^{}j^{}}`$ to be the number $`ii^{}`$. We say an empty square of $`\lambda _{\overline{3}}/\lambda _2`$ or an unplaced $`1`$-token is ready if its readiness number is $`0`$. Tokens which have reached their final destination and non-empty squares are not considered ready.

The key properties of this numbering scheme are the following:

###### Lemma 3.1.

Initially, the readiness numbers of the unplaced $`1`$-tokens are

1. weakly increasing along each row, and
2. weakly decreasing down each column;

in $`\lambda _{\overline{3}}/\lambda _2`$ the readiness numbers of empty squares are

1. weakly decreasing along each row, and
2. weakly increasing down each column.

Moreover, the GRGA, described below, preserves all of these properties.

###### Proof.

(a) If $`A`$ and $`B`$ are squares of $`\lambda _1`$ in the same row, and $`A`$ is right of $`B`$, then by Definition 3.2 $`f(A)`$ is lexicographically before $`f(B)`$. In particular, $`f(A)`$ is weakly above $`f(B)`$, i.e. in the same row or a row above. Thus the readiness number of the token in $`B`$ $``$ the readiness number of the token in $`A`$.

(b) If $`A`$ and $`B`$ are squares of $`\lambda _1`$, and $`A`$ is one square above $`B`$, then again $`f(A)`$ is lexicographically before $`f(B)`$. There are two cases. If $`f(A)`$ is strictly above $`f(B)`$, then the readiness number of the token in $`B`$ $``$ the readiness number of the token in $`A`$. Otherwise, $`f(A)`$ and $`f(B)`$ are in the same row in $`\lambda _{\overline{3}}/\lambda _2`$, with $`f(B)`$ right of $`f(A)`$. But then by Definition 3.2, $`B`$ must be lexicographically before $`A`$, which is a contradiction.

Statements (c) and (d) are proved similarly. That the GRGA preserves all these properties will be quite evident. ∎

Note that since the readiness number of the lower-leftmost token is at least $`0`$, by Lemma 3.1 parts (a) and (b) the readiness numbers are initially all non-negative.

#### The Grassmannian root game algorithm (GRGA).

The algorithm assumes that the tokens are in the initial positions of the root game for $`\pi _1,\pi _2,\pi _3`$ (corresponding to $`\sigma ,\mu ,\nu `$ with $`N`$ suitably large), that all $`3`$-tokens have been split into their own one-square region, and that we have a picture $`f`$ between $`\lambda _1`$ and $`\lambda _{\overline{3}}/\lambda _2`$. All moves take place in the big region.

1. If any of the $`1`$-tokens are ready, go to Step 2. Otherwise, perform a sequence of moves to shift all unplaced $`1`$-tokens up one square. The reader can easily check that the sequence of moves $`(1,2),(2,3),\mathrm{},(nl1,nl)`$ accomplishes this. The assumption that $`N`$ is sufficiently large ensures that the upward movement of the $`1`$-tokens is unobstructed. This step will decrease the readiness number of each $`1`$-token by $`1`$. Repeat this step until some $`1`$-token is ready.
2. Scan through the columns of $`\lambda _{\overline{3}}/\lambda _2`$, beginning with the rightmost column and proceeding to the left. Within each column locate the topmost square that does not already contain a $`1`$-token. Let $`S`$ be the first ready square which we encounter in this way.
3. Find a ready token $`t`$ in the same row as the square $`S`$. Make the unique move which causes $`t`$ to move into $`S`$. This may cause other tokens to move as well. All tokens which move here move to their final destination, so after this move, they and the squares they occupy are no longer considered ready.
4. Repeat Steps 1 through 3 until every square of $`\lambda _{\overline{3}}/\lambda _2`$ contains a $`1`$-token.

###### Example 3.4.

Figures 3 and 4 illustrate the GRGA, applied to the example from Figure 2. We now draw only the squares in the upper right $`3\times 7`$ rectangle as these are the only ones relevant to the movement of the $`1`$-tokens. Moreover, only the $`1`$-tokens are shown in these diagrams, and the number on the token is the readiness number, not the token label. To specify the picture $`f`$, each $`1`$-token is given a shading and the corresponding square under $`f`$ in $`\lambda _{\overline{3}}/\lambda _2`$ is shaded similarly. Each unshaded square actually contains a $`2`$-token. The two darkly shaded squares in the upper right corner contain $`3`$-tokens, as do each of the squares not shown in this diagram, but these squares are not part of the big region.

We now show that the GRGA accomplishes what it claims to accomplish.

###### Lemma 3.2.

Given a picture $`f`$ between $`\lambda _1`$ and $`\lambda _{\overline{3}}/\lambda _2`$, the GRGA will win the game for $`\pi _1,\pi _2,\pi _3`$.

###### Proof.

First note that after vertical movement from Step 1 is finished, the number of ready tokens in any row equals the number of ready squares in that row: an empty square $`S`$ is ready if and only if the $`1`$-token which began in the square $`f^1(S)`$ is in the same row as $`S`$. We show that this equality is preserved, by showing that the move in Step 3 only ever causes ready tokens to move into ready squares.

#### Claim (i)

A move from Step 3 causes only ready tokens to move. In particular the number of ready tokens in a row always remains less than or equal to the number of ready squares in a row.  The only tokens that can conceivably move are those in the same column as $`t`$. Because $`S`$ is the top unfilled square in its column, no tokens above $`t`$ move. Because of Lemma 3.1 part (b), all tokens below $`t`$ are ready.

#### Claim (ii)

Only ready squares are filled. The algorithm attempts to fill the rightmost squares first. If there is a ready square $`S^{}`$ in some column, the topmost empty square in that column will also be ready, by Lemma 3.1 part (d); thus the algorithm will never fill any square left of $`S^{}`$ before it fills $`S^{}`$. However, by Lemma 3.1 part (c) the ready squares are rightmost in their row. Thus if a token moves into a non-ready square, it means that there are no ready squares in its row. But since only ready tokens move, we would have a row with at least one ready token and no ready squares. This, as noted in Claim (i), is impossible.

#### Claim (iii)

The move from Step 3 is always possible. Since only ready squares are filled by ready tokens, the number of ready squares and ready tokens in any given row is always equal. Thus there is a ready token $`t`$ in the same row as the ready square $`S`$. Because we assume $`N`$ to be sufficiently large, $`t`$ is to the left of $`S`$. If $`t`$ is in column $`i`$ and $`S`$ is in column $`j`$, the move $`(i,j)`$ will take the token $`t`$ into square $`S`$.

Thus in Steps 2 and 3, every ready square eventually gets filled by a ready token: by the argument in Claim (ii) no square is skipped. However, because the readiness numbers are initially non-negative and Step 1 decreases the readiness number of each square by $`1`$, every square of $`\lambda _{\overline{3}}/\lambda _2`$ is ready at some point; thus the algorithm puts a $`1`$-token in each square of $`\lambda _{\overline{3}}/\lambda _2`$, at which point the game is won. ∎

###### Proof of Theorem 2.

($``$) This follows from Proposition 2.1 and Theorem 1.

($``$) For $`s=1`$ we use the GRGA, which wins the game by Lemma 3.2. For $`s>1`$, we proceed by induction. Suppose $`_{Fl(n+N)}[X_{\pi _1}]\mathrm{}[X_{\pi _{s+2}}]0`$. Then we can write

$$[X_{\pi _2}]\mathrm{}[X_{\pi _{s+1}}]=c[X_\rho ]+\mathrm{}$$
(5)
in the Schubert basis, where $`c>0`$, and

$$_{Fl(n+N)}[X_{\pi _1}][X_\rho ][X_{\pi _{s+2}}]0.$$

Since this is really a Grassmannian calculation, $`\rho `$ will be necessarily be of the form $`\pi ^{}(\sigma ^{},N)`$ for some $`0^{nl}1^l`$-string $`\sigma ^{}`$. By Lemma 3.2 we can win the game corresponding to $`\pi _1,\rho ,\pi _{s+2}`$, only moving $`1`$-tokens. It is easy to see that exactly the same sequence of splittings and moves can be made in the game for $`\pi _1,\pi _2,\mathrm{},\pi _{s+1},\pi _{s+2}`$, and that it causes the $`1`$-tokens to end up in exactly the same final positions. Note that we end up with either a $`1`$-token or an $`(s+2)`$-token in every square which does not correspond to an inversion of $`\rho `$, i.e. every square which does correspond to an inversion of $`w_0\rho `$.

This sequence of moves no longer wins the game; however, we can proceed inductively, after two further small steps. First, we perform a sequence of splittings so that every $`1`$-token is in a one-square region of its own. Next, after splitting in this way, we replace each $`1`$-token by an $`(s+2)`$-token. This second step is not a legitimate play in the game, but it is completely harmless: because every $`1`$-token is sequestered in its own one-square region, it can have no effect whatsoever on any possible subsequent moves of the game. But now we have precisely reached the initial position of the game corresponding to $`\pi _2,\mathrm{},\pi _{s+1},w_0\rho `$. This is again a game associated to a Grassmannian problem, and by Equation (5) the Schubert intersection number is non-zero. By induction, there is a sequence of moves to win this new game. Thus by concatenating the two sequences of moves, we can win the original game. ∎

## 4 Remarks

In Step 3 of the GRGA, there is a somewhat canonical choice for the token $`t`$, namely the leftmost ready token in its row. If we use this choice of $`t`$, one can verify that the algorithm actually transports the $`1`$-token which is initially in square $`S`$ to the square $`f(S)`$. On the other hand, if there is a way of winning the root game, there is generally a plethora of ways, most of which do not arise by following the GRGA for any picture. Theorem 2 tells us that the existence of any one way to win implies the existence of a picture between $`\lambda _1`$, and $`\lambda _{\overline{3}}/\lambda _2`$. However, given a sequence of moves which wins the game, it is not at all obvious how to construct such a picture. This is even unclear in the the case where $`s=1`$, and only $`1`$-tokens are moved.

It is worth noting that the root game can be used to determine whether

$$[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]0$$

even if the cohomological degree of the product is not $`dim_{}Gr_l(n)`$. To do this, we modify the game by changing the winning condition to read “the game is won if there is at most one token in each square”, rather than “exactly one token in each square”. Once we do this, we have the following corollary of Theorem 2.

###### Corollary 4.1.

Take $`N`$ suitably large, and let $`\pi _1,\mathrm{},\pi _{s+2}`$ be obtained from $`\sigma _1,\mathrm{},\sigma _s,\mu ,\nu `$ as before. Then $`[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]0`$ if and only if the root game for $`\pi _1,\mathrm{},\pi _{s+2}`$ can be won with the modified winning condition.

###### Proof.

Assume $`[Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]0`$. Then there exists $`\sigma ^{}`$ such that

$$_{Gr_l(n)}[Y_\sigma ^{}][Y_{\sigma _1}]\mathrm{}[Y_{\sigma _s}][Y_\mu ][Y_\nu ]0.$$

Let $`\pi ^{}=\pi (\sigma ^{},N)`$. Since we can win the unmodified game for $`\pi ^{},\pi _1,\mathrm{},\pi _{s+2}`$, we can win the modified game for $`\pi _1,\mathrm{},\pi _{s+2}`$ simply by omitting moves where the token corresponds to $`\pi ^{}`$. The reverse direction follows from Proposition 2.1 and \[P, Theorem 5\] (which generalises Theorem 1). ∎

There is a small catch: with this more general winning condition, our observation in Remark 1.1 becomes invalid. There is no longer an easy necessary and sufficient condition indicating when splitting is advantageous.

One of the unfortunate features of this presentation is the asymmetry in the way the permutations $`\pi _1,\mathrm{},\pi _{s+2}`$ are defined. The root game itself is manifestly symmetrical in the permutations given. However, because $`\pi _{s+1}`$ and $`\pi _{s+2}`$ are produced in a different way from $`\pi _1,\mathrm{},\pi _s`$, the symmetry is broken for Grassmannians. Nevertheless, as Theorem 2 is valid for any $`s`$, we can formulate a symmetrical game by taking $`\sigma _1,\mathrm{},\sigma _s`$ to be arbitrary, and $`\nu =\mu =0\mathrm{}01\mathrm{}1`$, so that $`[Y_\nu ]=[Y_\mu ]=1H^{}(Gr_l(n))`$. To see how this new formulation changes the initial position, contrast Figure 5 with Figure 2.

The only splitting which occurs in the proof of Theorem 2 is before the first move, and in the inductive step. The GRGA itself does not split between moves. As noted in Remark 1.1, it can never be harmful to split maximally between moves, and it turns out that if one modifies the GRGA to split maximally between moves, things proceed very much as before. However, in the next section, our proof of Theorem 3 will rely on the fact that the GRGA involves no splitting.

It would be nice if we could take $`N=0`$ in Theorem 2. Although we are not aware of any example which proves that this cannot be done, the algorithm simply falls apart if $`N`$ is too small. There are several problems which occur with trying to follow a similar approach. The most serious of these is that a token may be to the right of the square for which it is supposedly destined according to the chosen picture. Again, this highlights the fact that we do not know a straightforward two-way correspondence between pictures and ways of winning the root-game. Instead, in the next section we prove Theorem 3, a geometrical analogue of Theorem 2 which is valid for all $`N`$. Theorem 3 suggests that it is not unreasonable to conjecture that Theorem 2 is true for all $`N`$. We leave it as an open problem to determine whether or not this is in fact the case.

## 5 Geometric interpretation

In \[P\] we give a complete description of the geometry underlying the root game. The picture is quite different from those found in the geometric Littlewood-Richardson rules of Vakil \[V\] and more recently Coskun \[C\], both of which study degenerations of intersections of Schubert varieties inside a Grassmannian—we would be surprised if there were any straightforward relationship. Our methods are based on studying tangent spaces to Schubert varieties, and are more closely related to the approach used by Belkale in his geometric proof of the Horn conjecture \[B\]. Here, we will recall only the parts of the picture which are most relevant to our situation.

In this section we shall once again assume $`s=1`$. Our notation changes slightly from Section 1 in that we will be working with $`GL(n+N)`$ instead of $`GL(n)`$.

Let $`x_1,\mathrm{},x_{n+N}`$ denote the standard basis of $`^{n+N}`$. Let $`B`$ denote the standard Borel subgroup of $`GL(n+N)`$ (upper triangular matrices), and let $`B_{}`$ denote its opposite (lower triangular matrices). As before $`T=BB_{}`$ will be the standard maximal torus. For any complex vector space $`V`$ let $`Gr(V)`$ be the disjoint union of all Grassmannians $`Gr_d(V)`$, $`0ddimV`$.

Let $`R`$ be a region in the game, and let

$$\tau :R\text{subsets of }\{1,2,3\}$$

describe the configuration of the tokens within this region. In the underlying geometry there is, assigned to the combinatorial pair $`(R,\tau )`$, a corresponding geometric pair $`(V,U)`$, where $`V`$ is a $`B`$-module, and $`U=(U_1,U_2,U_3)`$ is a $`T^3`$-fixed point on $`Gr(V)^3`$ (or equivalently the $`U_k`$ are $`T`$-invariant subspaces of $`V`$). As a $`T`$-representation, $`V`$ is multiplicity-free, and the distinct $`T`$-weights are

$$\text{weights}(V)=\{x_ix_j|S_{ij}R\}.$$

The (distinct) $`T`$-weights of $`U_k`$ correspond to the positions of the $`k`$-tokens inside $`R`$:

$$\text{weights}(U_k)=\{x_ix_j|kS_{ij}\}.$$

Thus the pair $`(V,U)`$ carries all relevant information about the region $`R`$ and the arrangement of the tokens with $`R`$.

The region $`R`$ is solved when there is exactly one token in each square. In terms of the pair $`(V,U)`$ this is occurs when

$$V=U_1U_2U_3.$$
(6)
We’ll call any $`U=(U_1,U_2,U_3)`$ which satisfies condition (6) transverse.

Assuming we do not split the region $`R`$, a move or a sequence of moves in the game takes the pair $`(V,U)`$ to a new pair $`(V,U^{})`$, where $`U^{}`$ is in the $`B^3`$-orbit closure through $`UGr(V)^3`$. Thus if we solve a region starting from position $`(V,U)`$, we have located a transverse $`T^3`$-fixed point $`U^{}\overline{B^3U}`$.

The importance of transverse points in $`\overline{B^3U}`$ is seen in the following proposition.

###### Proposition 5.1 (\[P\]).

Consider the position of the root game game for $`\pi _1,\pi _2,\pi _3`$ which arises after splitting but before the first move is made. The tokens are in their initial position, but there may be more than one region. To each region $`R`$ there is an associated pair $`(V_R,U_R)`$. Then

$$_{Fl(n)}[X_{\pi _1}][X_{\pi _2}][X_{\pi _3}]1$$

if and only if for every region $`R`$ there exists a transverse point $`U^{}\overline{B^3U_R}`$.

Note that the point $`U^{}`$ in Proposition 5.1 is not necessarily $`T^3`$-fixed. The big question, therefore, is how specialised can we make the point $`U^{}`$ and still have Proposition 5.1 be true. There are three levels of specialisation that we could request of this transverse point $`U^{}`$.

1. $`U^{}`$ is any transverse point in $`\overline{B^3U_R}`$.
2. $`U^{}`$ is a $`T^3`$-fixed transverse point in $`\overline{B^3U_R}`$.
3. $`U^{}`$ is a ($`T^3`$-fixed) transverse point in $`\overline{B^3U_R}`$, where $`(V_R,U^{})`$ comes from applying sequence of root game moves (but no splitting) to the position $`(V_R,U_R)`$.

A priori, it is not clear that these three levels of specialisation are equivalent. However, for Grassmannian Schubert calculus with $`N`$ sufficiently large, Theorem 2 shows that they are all equivalent. The GRGA tells us exactly how to produce a sequence of moves which gives the point $`U^{}`$ at Level 3, which is the most specialised. Note it is important here that splitting is never used in the GRGA—when a region is split, the $`U^{}`$ one is tempted to define need not be in $`\overline{B^3U_R}`$.

Unfortunately, when $`N`$ is small the GRGA can fail, and so we cannot claim that all three levels of specialisation are equivalent for all $`N`$. Our goal in this section is to show that even if $`N`$ is too small for the GRGA to work, we can still get $`U^{}`$ at Level 2; i.e. Proposition 5.1 is still true for Grassmannian Schubert calculus if we demand that $`U^{}`$ be a $`T^3`$-fixed point.

To make matters more concrete, we now explicitly describe the initial pair $`(V,U)`$ for the big region $`R`$ in the root game associated to $`\sigma ,\mu ,\nu `$. This is the only region that we need to concern ourselves with, since all other regions are already solved.

Let $`M(n+N)`$ be the space of $`(n+N)\times (n+N)`$ matrices, having standard basis $`\{\stackrel{~}{e}_{ij}\}`$, and let $`B`$ act on $`M(n+N)`$ by conjugation. Let $`W`$ denote the $`B`$-submodule of $`M(n+N)`$ generated by the entries in the upper right $`(nl)\times (l+N)`$ rectangle. Let $`W^{}`$ be the $`B`$-submodule of $`W`$ generated by $`\stackrel{~}{e}_{ij}`$ such that $`S_{ij}`$ contains a $`3`$-token. Then $`V`$ is the quotient $`B`$-module

$$V=W/W^{}.$$

Note $`V`$ has a basis $`\{e_{ij}:=\stackrel{~}{e}_{ij}+W^{}|S_{ij}R\}`$. The point $`UGr(V)^3`$ is described as follows:

$$U_k=\text{span}\{e_{ij}|kS_{ij}\}.$$

Note that $`U_3=\{0\}`$, so we need not give it much further consideration.

Let $`U_{\overline{23}}`$ be the subspace of $`V`$ whose $`T`$-weights correspond to $`\lambda _{\overline{3}}/\lambda _2`$:

$$U_{\overline{23}}=\text{span}\{e_{ij}|2S_{ij}\}$$

and let $`U^{}=(U_{\overline{23}},U_2,U_3)`$. Note that $`V=U_{\overline{23}}U_2U_3`$, so the point $`U^{}Gr(V)^3`$ is transverse.

###### Example 5.1.

For the initial position shown in Figure 2,

$`V`$ $`=\text{span}\{\stackrel{~}{e}_{ij}|1i3,4j10\}/\text{span}\{\stackrel{~}{e}_{110},\stackrel{~}{e}_{210}\}`$
$`=\text{span}\{e_{14},e_{15},\mathrm{},e_{19},e_{24},e_{25},\mathrm{},e_{29},e_{34},e_{35},\mathrm{},e_{39},e_{310}\},`$
$`U_1`$ $`=\text{span}\{e_{14},e_{24},e_{25},e_{34},e_{35},e_{36}\},`$
$`U_2`$ $`=\text{span}\{e_{14},e_{15},e_{16},e_{17},e_{24},e_{25},e_{26},e_{27},e_{34},e_{35},e_{36},e_{37},e_{38}\},`$
$`U_{\overline{23}}`$ $`=\text{span}\{e_{18},e_{19},e_{28},e_{29},e_{39},e_{310}\}.`$
###### Theorem 3.

For every $`N0`$,

$$_{Gr_l(n)}[Y_\sigma ][Y_\mu ][Y_\nu ]1$$

if and only if with $`(V,U)`$ and $`U_{\overline{23}}`$ as above, $`U_{\overline{23}}\overline{BU_1}`$.

###### Proof.

$`()`$ If $`U_{\overline{23}}\overline{BU_1}`$ then $`U^{}`$ is transverse point in $`\overline{B^3U}`$. Hence this follows from Propositions 2.1 and 5.1.

$`()`$ Assume $`_{Gr_l(n)}[Y_\sigma ][Y_\mu ][Y_\nu ]1.`$

We know the result is true for $`N`$ sufficiently large, since the GRGA tells us how to get from the position $`(V,U)`$ to the position $`(V,U^{})`$. We use this fact to deduce the result for other values of $`N`$.

For any two choices of $`N`$, say $`N^\mathrm{a}`$ and $`N^\mathrm{b}`$, we get different spaces $`V`$, $`U`$, etc. We distinguish these notationally by using $`V^\mathrm{a}`$ (resp. $`V^\mathrm{b}`$) to denote the space $`V`$ corresponding to $`N=N^\mathrm{a}`$ (resp. $`N^\mathrm{b}`$), and likewise for any quantity depending on $`N`$. Note that $`d=dimU_1`$ is independent of $`N`$.

For any fixed $`N^\mathrm{a}`$ and $`N^\mathrm{b}`$, let $`\varphi :V^\mathrm{a}V^\mathrm{b}`$ be the linear map given by

$$\varphi (e_{ij})=\{\begin{array}{cc}e_{ij^{}}\hfill & \text{where }j^{}=j+N^\mathrm{b}N^\mathrm{a}\text{}S_{ij^{}}R^\mathrm{b}\hfill \\ 0\hfill & \text{if }S_{ij^{}}R^\mathrm{b}.\hfill \end{array}$$

Let $`A`$ be the dense open subset of $`Gr_d(V^\mathrm{a})`$,

$$A=\left\{XGr_d(V^\mathrm{a})\right|X\mathrm{ker}\varphi =\{0\}\}.$$

Then $`\varphi `$ induces a map $`\varphi _{}:AGr_d(V^\mathrm{b})`$

$$\varphi _{}(X)=\mathrm{Image}\varphi |_X.$$

Observe that $`\varphi _{}(U_{\overline{23}}^\mathrm{a})=U_{\overline{23}}^\mathrm{b}`$. The idea is essentially to show that $`AB^\mathrm{a}U_1^\mathrm{a}`$ is dense in $`\overline{B^\mathrm{a}U_1^\mathrm{a}}`$, and that $`\varphi _{}(AB^\mathrm{a}U_1^\mathrm{a})\overline{B^\mathrm{b}U_1^\mathrm{b}}`$. This implies that if $`U_{\overline{23}}^\mathrm{a}\overline{B^\mathrm{a}U_1^a}`$, then $`U_{\overline{23}}^\mathrm{b}=\varphi _{}(U_{\overline{23}}^\mathrm{a})\overline{B^\mathrm{b}U_1^b}`$. Hence if the result is true for $`N=N^\mathrm{a}`$, then the result will be true for $`N=N^\mathrm{b}`$.

We will only prove this in the case where $`N^\mathrm{a}=0`$ and $`N^\mathrm{b}`$ is arbitrary, and in the case where $`N^\mathrm{a}`$ is arbitrary and $`N^\mathrm{b}=0`$. This is enough to give the result for all $`N`$.

The case where $`N^\mathrm{a}=0`$ is the easier of the two. The map

$$\left[\begin{array}{cc}A^{nl\times nl}& B^{nl\times l}\\ 0& C^{l\times l}\end{array}\right]B^\mathrm{a}\left[\begin{array}{ccc}A^{nl\times nl}& 0& B^{nl\times l}\\ 0& I^{N^\mathrm{b}\times N^\mathrm{b}}& 0\\ 0& 0& C^{l\times l}\end{array}\right]B^\mathrm{b}$$

allows us to view $`B^\mathrm{a}`$ as a subgroup of $`B^\mathrm{b}`$, and $`\varphi _{}`$ is a $`B^\mathrm{a}`$-equivariant inclusion. Moreover $`\varphi _{}(U_1^\mathrm{a})\overline{B^\mathrm{b}U_1^\mathrm{b}}`$. Thus $`\varphi _{}`$ takes $`B^\mathrm{a}U_1^\mathrm{a}`$ into $`\overline{B^\mathrm{b}U_1^\mathrm{b}}`$.

For the case where $`N^\mathrm{b}=0`$, we consider the $`B`$-orbit not through $`U_1Gr_d(V)`$, but through a lifted point $`\stackrel{~}{U}_1Gr_d(W)`$. $`\stackrel{~}{U}_1`$ is defined in the same way as $`U_1`$: $`\stackrel{~}{U}_1=\text{span}\{\stackrel{~}{e}_{ij}|S_{ij}\text{ contains a }1\text{-token}\}`$. Let $`\stackrel{~}{\varphi }_{}`$ be defined analogously to $`\varphi _{}`$, taking a dense subset of $`Gr_d(W^\mathrm{a})`$ to $`Gr_d(W^\mathrm{b})`$. It suffices to show that $`\stackrel{~}{\varphi }_{}`$ takes a dense subset of $`B^\mathrm{a}\stackrel{~}{U}_1^\mathrm{a}`$ to a subset of $`B^\mathrm{b}\stackrel{~}{U}_1^\mathrm{b}`$.

Let $`LGL(nl)\times GL(l+N)`$ be the subgroup of $`GL(n+N)`$ of block diagonal matrices of type $`(nl,l+N)`$. Now $`L`$ also acts on $`W`$, and $`\stackrel{~}{U}_1`$ is fixed by $`B_{}L`$. Since $`(BL)(B_{}L)`$ is dense in $`L`$, it follows that the orbit $`B\stackrel{~}{U}_1=(BL)\stackrel{~}{U}_1`$ is dense in $`L\stackrel{~}{U}_1`$.

Thus in fact it suffices to show that $`\varphi _{}`$ takes a dense subset of $`L^\mathrm{a}\stackrel{~}{U}_1^\mathrm{a}`$ to a subset of $`L^\mathrm{b}\stackrel{~}{U}_1^\mathrm{b}`$. But this is true, as

$$\varphi _{}\left(\left[\begin{array}{ccc}A^{nl\times nl}& 0& 0\\ 0& B^{l\times N^\mathrm{a}}& C^{l\times l}\\ 0& D^{N^\mathrm{a}\times N^\mathrm{a}}& E^{N^\mathrm{a}\times l}\end{array}\right]\stackrel{~}{U}_1^\mathrm{a}\right)=\left[\begin{array}{cc}A^{nl\times nl}& 0\\ 0& C^{l\times l}\end{array}\right]\stackrel{~}{U}_1^\mathrm{b}$$

whenever both matrices are invertible.

Although Theorem 3 is a geometric statement, our proof ultimately relies on the combinatorics of the Littlewood-Richardson rule. The key non-geometric fact we use is that $`_{Gr_l(n)}[Y_\sigma ][Y_\mu ][Y_\nu ]0`$ if and only if there exists a picture between $`\lambda _1`$ and $`\lambda _{\overline{3}}/\lambda _2`$. It would be an interesting project to find a purely geometric proof of this theorem. The hope would be that a geometric proof of Theorem 3 might allow us to see how to generalise some of the results in this paper beyond the Grassmannian.

### Acknowledgements

The author is grateful to Allen Knutson and Stephanie van Willigenburg for providing feedback and corrections on this paper.


text:
# Correlations between pressure and bandwidth effects in metal-insulator transitions in manganites

## Abstract

The effect of pressure on the metal-insulator transition in manganites with a broad range of bandwidths is investigated. A critical pressure is found at which the metal-insulator transition temperature, TMI, reaches a maximum value in every sample studied. The origin of this universal pressure and the relation between the pressure effect and the bandwidth on the metal-insulator transition are discussed.

Manganites have been the focus of intense studies in recent years since the observation of large magnetoresistance sparked interest in these materials for use as magnetoresistance sensors. The metal-insulator transitions (MIT) observed in this kind of materials is crucial to the colossal magnetoresistance effect. Metal-insulator transition occurs in manganites in two cases: first, metallic ground state exists in the low temperature range in some doping systems at certain doping concentrations, such as: La1-xSrxMnO3 (x$``$0.16-0.50), La1-xCaxMnO3 (x$``$0.18-0.50), Nd1-xSrxMnO3 (x$``$0.25-0.50); second, metallic states can be induced by other factors, such as magnetic fields, photons, pressure, electric fields. Pr1-xCaxMnO3 at x$``$0.3 is typical in the latter class. In most of the manganites, the metallic state is coupled to the ferromagnetic state so that the very large magnetoresistance can be explained by double exchange theory.

In the parameters determining the complicated properties of the manganites, the eg electron bandwidth W is a particularly important one to the metal-transition temperature TMI, or the appearance of the metallic state under some factors. In manganites, the Mn 3d orbital is split into t2g and eg orbitals by the octahedral crystal field. The conduction band electrons are of eg symmetry. Because the eg orbital is Jahn-Teller active, Jahn-Teller distortion (JTD) can further split the two-fold degenerate eg orbital to trap the conduction band electrons. Consequently, the bandwidth is highly correlated with the local atomic structure of the MnO6 octahedra: cooperative tilting (Mn-O-Mn bond angle), Jahn-Teller distortion (Mn-O distances) and coherence of the JTD. The bandwidth is characterized by the overlap between the Mn-3d orbital and O-2p orbital and can be described empirically by an equation:Medarde et al. (1995)

$$W\frac{\mathrm{cos}[\frac{1}{2}(\pi <\beta >)]}{d_{MnO}^{3.5}}$$
(1)
where W is the bandwidth, $`\beta `$ is the Mn-O-Mn bond angle, and dMn-O is the Mn-O bond length. In double exchange theory, it is described as the electron hopping rate or the transfer integral: $`t_{ij}=t_{ij}^0\mathrm{cos}(\frac{\theta _{ij}}{2})`$, where $`t_{ij}^0`$ is the transfer integral that depends on the spatial wave function overlaps, $`\theta _{ij}`$ is the relative angle between two neighboring Mn ion t2g core spins.

Generally, the structure can be tuned in two ways: chemical doping and external pressure. In chemical doping, by selecting different doping elements and doping concentration, the average A-site atom size $`<`$rA$`>`$ in the AMn$`O_3`$ system is changed. Because of the mismatch between $`<`$rA$`>`$ and Mn-site ion size, the local atomic structure of MnO6 octahedra can be modified. Therefore, the bandwidth is tuned by chemical doping so that complicated electronic and magnetic phase diagrams have been observed.Tokura and Tomioka (1999) The external pressure method, is a “clean method” that only modifies lattice structure without inducing chemical complexity. To date, in studies on manganites, the effects of external pressure on the charge ordering, metal-insulator transition, magnetic states have been observed.

Currently, most of the high pressure studies on manganites are on metal-insulator transition. In the low pressure range, this electronic transition is coupled to the ferromagnetic transition, which can be explained qualitatively by double exchange theory. Neumeier et al. (1995); Moritomo et al. (1997) It is also found that hydrostatic pressure has a similar effects to chemical doping with larger atoms and higher doping concentration. Both can increase the Mn-O-Mn bond angle, compress Mn-O bond length and hence, lead to larger bandwidth. Correspondingly, TC (or TMI) increases, or in some manganites originally in insulating state, a MIT is induced. The effect of chemical doping and pressure can be scaled to each other with a conversion factor 3.75$`\times `$10-4 Å/kbar.Hwang et al. (1995) However, most pressure experiments were conducted below 2 GPa.

By applying pressures up to $``$6 GPa on manganite systems with a broad range of bandwidths, the effect of pressure on the MIT and the correlation between the pressure effect and the chemical doping were observed. It is found that TC and/or TMI do not change monotonically with pressure Cui et al. (2003a) and these two transitions do not always couple.Cui and Tyson (2003) A universal pressure may exist for the metal-insulator transition in manganites. With an increase of the bandwidth, the change in the metal-insulator transition temperature with pressure may vanish.

To systemically explore the external pressure effect and the chemical doping effect and the correlation between them, several manganite systems with much different ground magnetic and electronic states, Nd1-xSrxMnO3 (x = 0.45, 0.50) (NSMO), La0.60Y0.07Ca0.33MnO3 (LYCMO), and Pr1-xCaxMnO3 (x = 0.25, 0.30, 0.35) (PCMO), were selected. The bandwidth of the selected samples uniformly distributes in the electronic phase diagram (see Ref. Hwang et al., 1995).

The samples were prepared by solid-state reaction. The procedure and details of making the samples were described elsewhere.Cui et al. (2003a); Cui and Tyson ; Cui et al. (2003b) All the samples are characterized with the x-ray diffraction and magnetization measurements. The details of high pressure resistivity measurement method and error analysis were described previously.Cui et al. (2003a) The metal-insulator transition temperature whenever present is defined as the temperature at the resistivity peak. Because of the lower temperature stability of our system in the cooling cycle, the data were taken only while warming up.

In all the samples studied, there is a MIT at ambient pressure or a MIT can be induced by applying pressure. Corresponding to the bandwidth phase diagram in Ref. Hwang et al., 1995, the Nd1-xSrxMnO3 (x = 0.45, 0.50) system has a large bandwidth; La0.60Y0.07Ca0.33MnO3 has a medium bandwidth; Pr1-xCaxMnO3 (x = 0.25, 0.30, 0.35) system has a small bandwidth.

In Table 1, $`<`$rA$`>`$, the tolerance factor t, and metal-insulator transition temperature at ambient pressure, which corresponds to the bandwidth, are listed. The average Mn-O bond length and Mn-O-Mn bond angle of all samples determined from the Rietveld refinement to the x-ray diffraction patterns are shown in Fig. 1. According to equation (1), with increasing $`<`$rA$`>`$ or t, the decreasing bond length and bond angle lead to increasing bandwidth W and hence, increasing TMI.

With the application of pressure, the metal-insulator transition temperatures of the samples which have MIT at ambient pressure increases. In the narrow bandwidth Pr1-xCaxMnO3 system, the samples are insulating at ambient pressure. Under pressure, metal-insulator transitions are induced. With pressure increase, the behavior of the TMI is similar to other samples with larger bandwidth. When the pressure is above a certain point, the increasing trend of TMI of all samples is reversed. The evolution of the transition temperatures of all samples are shown in Fig. 2.

In Fig. 2, the most salient feature is that a critical pressure P\* exists in each sample: with pressure increase, below P\*, TMI increases; above P\*, TMI decreases. By fitting TMI vs. P plots with a third-order polynomial, P\* for each sample can be extracted and is listed in Table 1. In the fitting error range, the samples have the same critical pressure P\* (One exception is the NSMO45 sample, the P\* determined from the TMI is small. But if we look at the resistivity changes with pressure, its critical pressure is the same as other samples. This may comes from the highest temperature limit of the instruments which leads that the TMI near to the limit in the middle range pressure can not be determined.Cui et al. (2003b)).

In large bandwidth samples, the change of TMI with pressure is slower than in narrow bandwidth ones, indicating that the large bandwidth samples are more stable under pressure. The samples studied are selected with different doping concentration and from different doping systems. The bandwidths spans a large range. The samples also have much different ground state electronic and magnetic properties at ambient conditions. But the metal-insulator transitions in these samples all follow a similar behavior, therefore, it is reasonable to speculate that the critical pressure P\* is universal for the metal-insulator transitions in manganites. From the structural measurements on manganites,Cui et al. (2003a); Congeduti et al. (2001); Meneghini et al. (2002) the behavior of TMI under pressure could possibly be ascribed to a local atomic structure transformation of the MnO6 octahedra.

Under pressure, the smaller bandwidth samples seem to have smaller pressure range in which the sample is metallic in low temperature and outside which they are insulating. On the other hand, the samples with large bandwidth is more stable under pressure and the variable range of TMI is small and they do not become insulating in a larger pressure range. The lower stability of TMI in small bandwidth samples may come from the small A-site atoms, which leave more space between the octahedra for them to rotate - accordingly, a smaller pressure window for the metallic state.

In Fig. 2, the large bandwidth samples have higher TMI. The only exception is that in the Nd1-xSrxMnO3 system, the x = 0.5 compound nominally has a larger bandwidth than the x = 0.45 compound but has lower TMI. This possibly results from the strong charge ordering effect in Nd0.5Sr0.5MnO3.

The changing rate of the metal-insulator transition temperature with pressure at the ambient pressure, dTMI/dP at P = 0, is also interesting. The values of dTMI/dP extracted from the third-order polynomial fitting results are listed in Table 1. Clearly, the smaller the bandwidth, the larger dTMI/dP, indicating that the local structure of smaller bandwidth sample is more distorted and has a relatively large degree to which it can be compressed by pressure.

In summary, by applying external pressure on manganites of different chemical doping system and doping concentration and hence different eg electron bandwidth, it is found that the pressure effect on the metal-insulator transition in manganites is not equivalent to that of the chemical doping. Only at low pressures, is the pressure effect on the metal-insulator transition analogy to the chemical doping with elements with large atom size. With pressure increase, the trend of TMI increase with pressure is reversed at a critical pressure, above which the transition temperature decreases with pressure and finally, the material may become insulating. The critical pressure is found to exist in all the samples studied and possibly is universal for the metal-insulator transition in the manganites. The bandwidth (chemical doping) determines how stable the material may be under pressure. The larger bandwidth manganites are more stable under pressure and therefore, have smaller dTMI/dP near ambient pressure and smaller TMI variation under pressure. Because of the importance of the local atomic structure of the MnO6 octahedra to the electronic and magnetic properties of the manganites, this work may also contribute to understanding the properties of thin films of manganites which are important in technological applications.

This work is supported by National Science Foundation Grant DMR-0209243.


text:
# The Diverse Nature of Intrinsic Absorbers in AGNs

## 1. Introduction: Why Study Intrinsic Absorbers?

Let us start by defining “intrinsic” absorption in terms of gas that is (or was) part of the overall AGN/host galaxy environment. This definition excludes only very distant, cosmologically intervening material, such as intergalactic clouds or unrelated galaxies. It reminds us that, especially in quasar studies, absorption can occur in a wide range of environments. The rich variety of intrinsic absorbers yields numerous diagnostics of both the AGN phenomenon and the AGN–host galaxy connection. A short list of reasons for studying intrinsic absorption might include the following.

$``$ Intrinsic absorbers are a fundamental component of AGN environments. They are common in Type I (broad emission line) AGNs (see below), and might be ubiquitous if the absorbing gas fills only part of the sky as seen from the central continuum source. In addition, the amounts of absorbing gas might be enormous — rivalling or exceeding the mass in the broad emission line region.

$``$ Many intrinsic absorbers are involved in AGN outflows. The flows are driven by the same accretion processes that feed the central super-massive black hole (SMBH) and fuel other AGN energetics. The need for accreting matter to expel angular momentum probably means that the wind mass loss rates, $`\dot{M}_{wind}`$, are directly proportional to the mass accretion rate, $`\dot{M}_{acc}`$.

$``$ The relationship between outflow and accretion also implies that intrinsic absorbers are connected to the basic physics of SMBH growth and AGN evolution.

$``$ Intrinsic absorption that occurs far from the AGN might uniquely measure a variety of regions in the host galaxies, such as the interstellar medium, gas streams in galactic halos, or galactic super-winds driven by starburst activity.

$``$ The metal abundances in high-redshift intrinsic absorbers can provide unique constraints on the amount of star formation and the overall maturity of young galactic or proto-galactic nuclei.

$``$ The metal-rich gas expelled by high-redshift AGNs might be a significant source of metal “pollution” to the intergalactic medium at early cosmic times.

In this brief review, we focus on a few issues regarding absorption line classification, the relationships between classes, and the implications of partial coverage of the background light source(s). See also the reviews by Crenshaw, Kraemer, & George (2003), Hamann (2000), Hamann & Ferland (1999), and the ASP conference series volumes (128 and 255) devoted to AGN mass loss.

## 2. Absorption Line Classes

AGN absorption lines are classified empirically by the full width at half minimum (FWHM) of their profiles. The class definitions necessarily evolve as we encounter new phenomena and assimilate different measurement schemes. Nonetheless, standardized classes are essential. The main classes are the “narrow” absorption lines (NALs), the “broad” absorption lines (BALs), and a catch-all intermediate class called the mini-BALs.

BALs are blueshifted from the systemic (emission line) redshift by as much as $``$0.2$`c`$, and clearly form in winds form the central engine. Weymann et al. (1991) introduced a “BAL-nicity” index to define this class. Hall et al. (2002) proposed a less restrictive index to include a wider range of line widths. We strongly advocate the use of quantitative indices, but a reasonable starting point for casual conversation is that BALs have continuous absorption over $`>2000`$ km/s, with at least some portion of the profile having a velocity shift $`v>2000`$ km/s compared to systemic. BALs are further divided into subclasses according to the degree of ionization apparent in the lines. “HiBALs” have nominally SiIV $`\lambda \lambda `$1394,1403 and CIV $`\lambda \lambda `$1548,1551 (or perhaps CIII $`\lambda `$977) as their lowest ionization lines. “LoBALs” include lower ionization stages, such as MgII $`\lambda \lambda `$2796,2804. “FeLoBALs” have more extreme low-ionzation regions that produce excited-state absorption in FeII (requiring densities $`n_e>10^5`$ cm-3).

A useful (but physically arbitrary) definition of NALs is that they are narrow enough not to blend important UV doublets, e.g., the CIV pair with separation $``$500 km/s. Thus we require FWHM $`<200`$ to 300 km/s. NALs with velocity shifts $`v<5000`$ km/s from systemic are also call “associated” absorption lines (AALs) because of their plausible physical relationship to the AGN. No one has yet sub-divided AALs into HiAALs and LoAALs analogous to the BALs, but it would be interesting to compare sub-classes based on these properties.

Mini-BALs have FWHMs intermediate between the NALs and BALs. They appear at the same range of blueshifted velocities as the BALs, and they also clearly form in AGN outflows. Examples of high-velocity mini-BALs can be found in Jannuzi et al. (1996) and Hamann et al. (1997a). For the BALs, see Weymann et al. (1991), Reichard et al. (2003), and Hall et al. (2002). For the AALs, see Foltz et al. (1986) and Ganguly et al. (1999). See also the reviews cited above.

Figure 1 lists the major absorption line classes along with some properties and speculation that we will describe below. The detection frequencies given in the figure are our own educated guesses, except for the quasar AALs (G. Richards, priv. comm.), Seyfert 1 AALs (Crenshaw et al. 1999), and the HiBALs and LoBALs (Weymann et al. 1991, Reichard et al. 2003). Note that we list all Seyfert 1 absorbers as NALs, even though some may be considered mini-BALs.

## 3. Identifying Intrinsic NALs

Quasar NALs are known to form in a wide range of environments, from outflows like the BALs to unrelated gas or galaxies at large (cosmologically significant) distances from the AGN (Weymann et al. 1979). NALs in low-redshift Seyfert 1 galaxies might be almost exclusively intrinsic, but there is still a range of possible locations within the global AGN/host galaxy environment. An important goal, therefore, is to identify the intrinsic NALs and locate the absorbing regions with respect to the AGNs/host galaxies. One strategy is to examine the probabilities. For example, quasars have an over-density of NALs near the emission redshift, i.e., more than expected for a random distribution of cosmologically intervening material (Weymann et al. 1979, Foltz et al. 1986). Others have noted that the presence and/or strengths of NALs correlate with AGN properties such as luminosity, radio-loudness, and radio lobe orientation (Wills et al. 1995, Richards et al. 1999, Richards 2001). These relationships imply that a substantial fraction of quasar AALs ($`>`$50%), and probably some of the highly blueshifted ($`v>5000`$ km/s) NALs, are intrinsic to quasars.

We can also examine individual absorbers to look for i) absorption line variability, ii) profiles that are broad and smooth compared to thermal speeds, iii) partial line-of-sight coverage of the background light source, and iv) high gas densities based on excited-state absorption lines (see below, also Tripp, Lu, & Savage 1996, Hamann et al. 1997b, 2001, Barlow & Sargent 1997, Ganguly et al. 1999, Narayanan & Hamann 2003). These characteristics often appear together and suggest an intrinsic origin because they are most easily understood in terms of the dense and dynamic environments near AGNs.

## 4. Absorber Trends & Characteristics

Brandt, Laor, & Wills (2000) showed that the strengths of intrinsic absorption lines (e.g., $`W_\lambda `$(CIV)) correlate roughly with the amount of X-ray absorption. Moreover, the column densities derived for the X-ray absorbers are much larger than earlier estimates from the UV lines, to the point where some BAL quasars have Thompson thick X-ray absorbers with $`N_H>10^{24}`$ cm-2 (Hamann 1998, Gallagher et al. 1999, Mathur et al. 2000, Arav et al. 2002). In some Seyfert 1 galaxies (e.g., Kaspi et al. 2002) we know that the X-ray absorbing gas is outflowing at modest speeds with the UV line absorbers. It would be more problematic for models invoking radiative acceleration if the large X-ray columns in BAL quasars turn out to be outflowing at BAL speeds (Hamann 1998). In any case, we have an expectation that mass loss rates increase generally with line strength, as shown in Figure 1.

Other studies suggest that, in very rough terms, the amount of reddening, the percent polarization, and perhaps the IR luminosities are all nominally higher in BAL quasars and may increase toward lower ionizations through the BAL class (Weymann et al. 1991, Schmidt & Hines 1999, Brotherton et al. 2002, Hall et al. 2002, Richards et al. 2003, Reichard et al. 2003, and refs. therein). It is not known if these tendencies extend to the mini-BALs or AALs.

## 5. Evolution & Orientation

AGN evolution is something of a holy grail because we know almost nothing about it. Unlike stars and galaxies, we have yet to identify “young” or “old” AGNs. It has been suggested that the FeLoBAL and LoBAL quasars represent a young phase because their enormous column densities, low ionizations, and reddening might be signatures of an energetic AGN that is still partially enshrouded in its dusty parental environment (Voit, Weymann, & Korista 1993, Gregg et al. 2000, Brotherton et al. 2002, and refs. therein). Perhaps this evolutionary path continues through the BALs to the mini-BAL and AAL stages as the outflows gradually weaken (Figure 1). However, orientation probably also plays an important role. Most models have BAL winds rising up out of the accretion disk and flowing close the the disk plane (Murray et al. 1995, Proga, Stone, & Kallman 2000). Viewing angles close to the disk plane should plausibly lead to larger absorption column densities, detections of lower ionization stages, and perhaps increased reddening and polarization (Schmidt & Hines 1999, Elvis 2000). Similar arguments may apply to the AALs (Wills et al. 1995). It seems likely that orientation and evolution both contribute to differences between the absorber classes.

One observational difficulty is that AGN lifetimes are short compared to their host galaxies, so we cannot argue that young/old AGNs reside in young/old galaxies. In particular, AGN luminosities, $`L`$, are related to the mass accretion rate onto the central SMBH, $`\dot{M}_{acc}`$, by an efficiency factor, $`\eta `$. The luminosities are also some fraction, $`\gamma `$, of the theoretical Eddingtion limit, $`L_{edd}`$, such that

$$L\eta \dot{M}_{acc}c^2\gamma L_{edd}(1.5\times 10^{46})\gamma M_8(\mathrm{ergs}/\mathrm{s})$$
(1)
where $`M_8`$ is the SMBH mass in units of $`10^8`$ M. If we adopt standard values of $`\eta =0.1`$ and $`\gamma =0.5`$, and we assume the lifetime of the bright AGN phase is the time needed for the last doubling of the SMBH mass, then we derive a nominal lifetime of $``$$`6\times 10^7`$ yr (see also Ferrarese 2002 and refs. therein).

These lifetimes are also short compared to the duration of the main quasar epoch, for example, the $``$$`5\times 10^9`$ yrs elapsed from redshifts 4 to 1. Therefore, we cannot expect to find exclusively young/old AGNs at high/low redshifts. However, we can use the fact that the AGN population is growing rapidly at redshifts $`z>3`$. At these redshifts, there are more AGNs being created “now” than in the past, so there should be more young AGNs than old ones. Conversely, at $`z<1`$ the AGN population is in rapid decline and there should be more old than young AGNs. The strength of this young versus old signature depends mainly on the slope of the relation between AGN co-moving number density and time. The signature might be strong enough to study age-related trends with redshift in large databases such as the SDSS. We should also keep in mind that the truly youngest AGNs (preceding quasars and having low SMBH masses) will be found among the least luminous sources, perhaps as heavily enshrouded Seyfert 1-like nuclei in dusty starburst galaxies (Figure 1).

## 6. Column Densities & Partial Coverage

Column density is one of the most critical parameters in intrinsic absorber studies. If we have high resolution spectra, e.g., of absorption lines, we can derive the column density at each velocity shift, $`N_v`$, from the optical depth profile, $`\tau _v`$,

$$N_v=\frac{m_ec}{\pi e^2f\lambda }\tau _v$$
(2)
where $`\lambda `$ is the line wavelength and $`f`$ is the oscillator strength. In principle, this analysis is straightforward. In practice, there can be observational and other limitations. There is also a simple question we should consider, which is, what do we mean by “the column density” in a particular ion?

If you walk into a room full of astronomers, even AGN astronomers, and start talking about inhomogeneous absorbing media and partial coverage of the background light source, you will see some eyes roll back in their heads while others look longingly for relief at the clock on the wall. Nonetheless, we keep talking because we believe partial coverage is a fundamental problem that needs to be addressed right now in studies of intrinsic absorption. The problem is more common and probably more insidious than we previously believed. Continuing progress on the big goals outlined in §1 requires a better understanding of the partial coverage issue.

The simplest situation is homogeneous partial coverage (HPC), where the absorber has a constant optical depth across the projected area of an extended emission source. The essential geometry is illustrated in Figure 2a. For simplicity, we assume that the emission source has a constant intensity, $`I_o`$, across its projected area. If $`\tau _v`$ is the optical depth through the absorber at a given velocity shift, and $`C_f(v)`$ is the fraction of the emitting area covered by the absorber (with $`0C_f(v)1`$), then the measured intensity is a spatial average given by,

$$\frac{I(v)}{I_o(v)}=1C_f(v)+C_f(v)e^{\tau _v}$$
(3)
If we measure two lines with a known ratio of optical depths, as in a doublet, then we can solve for the two unknowns, $`C_f(v)`$ and $`\tau _v`$, using the two equations for the measured intensities (Barlow & Sargent 1997, Hamann et al. 1997b, Ganguly et al. 1999).

There is growing evidence that the reality is not like Figure 2a, but more like Figure 2b — where the absorber has a range of optical depths across the projected area of the emission source(s). We call this situation inhomogeneous partial coverage (IPC, see deKool et al. 2002, Sabra & Hamann 2003). The evidence for IPC absorbers comes from derived coverage fractions that depend on velocity and, at a given velocity, on the ionization and/or line strength (i.e., true optical depth, see below, also Barlow et al. 1997, Hamann et al. 2001, Arav et al. 2002, and refs. therein). In an IPC cloud distribution (Figure 2b), strong transitions in abundant ions can have $`\tau _v>1`$ over a larger area, leading to larger derived coverage fractions. Similarly, low ionization species might have lower coverage fractions if they are concentrated in the smaller, denser regions.

In general, there is a 2-dimensional optical depth structure, $`\tau _v(x,y)`$, in front of an extended emission source with intensities, $`I_o(x,y,v)`$, where $`x`$ and $`y`$ are the spatial coordinates. The observed intensity at each velocity shift is the $`I_o`$-weighted average of $`e^{\tau _v(x,y)}`$,

$$I(v)=I_o(x,y,v)e^{\tau _v(x,y)}\frac{dxdy}{A}$$
(4)
where $`A=𝑑x𝑑y`$ is the total projected area of the emission source (deKool et al. 2002, Sabra & Hamann 2003). deKool et al. (2002) constructed hypothetical $`\tau _v(x,y)`$ distributions and showed how they can affect measured line strengths. They also showed that any $`\tau _v(x,y)`$ can be represented by an equivalent 1-dimensional $`\tau _v(x)`$ without loss of generality. Sabra & Hamann (2003) extended the pioneering work of deKool et al. by i) considering a wider range of $`\tau _v(x)`$ possibilities, ii) showing that even in complex IPC situations we can use measured doublets (or multiplets) to estimate the $`\tau _v(x)`$ shape, and iii) comparing the results of IPC and HPC doublet analyses.

Figure 3 shows one hypothetical $`\tau _v(x,y)`$ distribution composed of overlapping “clouds,” each with a gaussian optical depth spatial profile (see Sabra & Hamann 2003). Also shown is the equivalent 1-dimensional $`\tau _v(x)`$. We use this optical depth distribution to calculate intensities for an absorption line doublet. As in the HPC analysis, the doublet intensities provide two equations from which we can derive two unknowns. But here we adopt a simple functional form for the optical depths, such as the power law $`\tau _v(x)=\tau _{max}x^a`$, and then solve for $`\tau _{max}`$ and the “shape” parameter $`a`$. The results are shown in Figure 3. The power law provides a close but imperfect match to the input $`\tau _v(x)`$ because the input was not a power law. To do better, e.g., in the analysis of real AGN lines, we need theoretical guidance on what functional forms of $`\tau _v(x)`$ are appropriate. It would also help to measure sets of $``$3 lines with known optical depth ratios, e.g., in multiplets, to obtain more constraints on $`\tau _v(x)`$.

Finally, we return to the question of what we mean by “the optical depth” or “the column density” in a given feature. Are there single values of these parameters that are meaningful in an IPC environment, e.g., that we can use in a simple ionization/abundance analysis? The answer is probably yes. The input $`\tau _v(x,y)`$ distribution in Figure 3 has a spatially averaged optical depth of $`\tau _{IPC}^{avg}1.2`$. The HPC doublet analysis applied to the line intensities emerging through that distribution yields $`\tau _{HPC}^{avg}1.1`$. Sabra & Hamann (2003) performed many such experiments. We found that $`\tau _{IPC}^{avg}/\tau _{HPC}^{avg}1`$ to within $``$30% or so, as long as the input $`\tau _v(x,y)`$ does not contain spatially narrow “spikes.” Such spikes (very large optical depths over a small coverage area) can profoundly affect the average optical depth while having little impact on the observed intensities. Therefore, even with potentially complex IPC absorbers, we can use spatially averaged optical depths derived from the HPC analysis to obtain spatially averaged estimates of the column densities, ionization, abundances, etc.

## 7. Applications to Real Data

We obtained (with our collaborators T. Barlow, E.M. Burbidge, & V. Junkkarinen) high-resolution spectra of a sample of AAL quasars using the HIRES spectrograph at the Keck 10 m telescope. The quasars have redshifts of order 2, allowing us to measure absorption lines in the rest-frame UV. Figure 4 shows line profiles for several doublets in two of these quasars. The HPC doublet analysis shows that these lines are optically thick at all velocities, while velocity-, ion-, and line-dependent coverage fractions determine the line strengths and profiles (see also Barlow et al. 1997, Arav et al. 2002, and refs. therein). Each individual line (or velocity component) probably does not correspond to a single entity we might call a “cloud” in the absorber. These features will provide only lower limits on the ionic column densities. This situation seems to be common for intrinsic quasar NALs with FWHMs $`>`$ 100 km/s. However, narrower quasar NALs appear to have typically $`C_f(v)1`$, with each distinct velocity component corresponding to a distinct cloud. This simpler circumstance can yield more reliable column densities and abundances (we find typically $`Z>`$ Z).

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