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# A Modified Theory Of Newtonian Mechanics

## Abstract

A specific model of the inertial law is presented by which we can have some deeper insight into the essence of mass and inertia. In this modified theory there is no need to keep the concept of absolute space and the third law as a principle. By introducing a convenient form for gravitational law the coupling constant G becomes a function of inertial parameters of the universe.

Keywords:Absolute Space,Classical Theories of Gravity,Cosmology,Inertia.

According to the Newton’s second law of motion the relation $`\stackrel{}{F}=m\stackrel{}{a}`$ holds between the applied force on a particle and its acceleration in absolute space and any other frames which are related to it by Galilean transformations. These reference frames are called inertal frames and $`m`$ is the inertial mass of the particle. The relation between coordinates of two inertial systems $`S`$ and $`S^{}`$ which are moving with constant velocity $`\stackrel{}{v}`$ with respect to each other are :

$$\{\begin{array}{c}t^{}=t\hfill \\ \stackrel{}{a}^{}=\stackrel{}{a}\hfill \\ \stackrel{}{u}^{}=\stackrel{}{u}\stackrel{}{v}\hfill \\ \stackrel{}{x}^{}=\stackrel{}{x}\stackrel{}{v}t+\stackrel{}{x}_0\hfill \end{array}$$
(1)
Despite its simple appearance and practical applications Newton’s concept of absolute space which is the basis of Newtonian mechanics (NM) and in turn the basis of mass and energy is not well defined and has been criticized. Cogent arguments against absolute space are :

* \- There is no unique way of locating Newton’s absolute space within the finite class of inertial frames.
* \- It conflicts with one’s understanding to conceive of a thing which acts but cannot be acted upon.

In his critique of Newton’s conceptions in “The Science of Mechanics”, Mach suggested a set of pioneering implicit ideas then addressed by A. Einstein as Mach’s Principle (MP) . According to MP, every motion is only comprehensive as a relative motion, inertia is due to an interaction with average mass of the universe  and the distribution of matter in the universe determines the inertial frame at each point. There have been many attempts to find Machianized reformulations of mechanics. Einstein’s main aim in the general theory of relativity (GR) was to do this through giving an equation for gravitation and inertia together. But it came out that his theory shows some non-Machian aspects and do not fulfill the MP. For instance it has solution for empty space. Also there are many alternative theories of gravitation motivated by Mach’s ideas . Since in GR as our standard theory of gravitation we are still faced with unresolved questions such as singularity problems and lack of well defined source for inertia, then it is reasonable to have more discussions about foundations. It seems with closer scrutiny of Newton’s second law we may find some insight.

Actually in this work our aim is to present a simple method for Machianization of NM. We propose that inertial effect is a mutual interaction between two particles which in any non-rotating arbitrary reference frame $`S`$ is proportional to the difference between their accelerations and to the inertial charges of each individual particles as follows

$$\stackrel{}{F}_{inertia}=\mu .c_1.c_2(\stackrel{}{a}_1\stackrel{}{a}_2)$$
(2)
where $`\stackrel{}{a}_1`$ , $`\stackrel{}{a}_2`$ are accelerations of particles 1 and 2 with respect to $`S`$ , their inertial charges are denoted by $`c_1`$ , $`c_2`$ respectively , and $`\mu `$ is an inertial coupling constant. This can be easily extended to systems consisting of N particles. Again in any non-rotating arbitrary reference frame $`S`$ we have

$$\stackrel{}{F}_i=\mu .c_i\underset{j=1}{\overset{N}{}}c_j(\stackrel{}{a}_i\stackrel{}{a}_j)$$
(3)
where $`\stackrel{}{F}_i`$ is the applied force on particle i and the summation is done over all particles.

In the real world the inertial charge and Newtonian inertial mass of a particle are related so that

$$m_i=\mu .c_i\underset{j=1}{\overset{all}{}}c_j$$
(4)
where summation is done over all particles in the universe. Since local inhomogenities have no observed effects on the inertial mass then it is accepted that the inertial mass is determined by the global structure of the universe and the relation (4) for the inertial mass is in accordance with this general belief. This can be accounted as a simple formulation of Mach’s idea about inertia. In terms of inertial mass , equation (3) can be rewritten in the following form

$$\stackrel{}{F}_i=m_i\left(\stackrel{}{a}_i\frac{\underset{j=1}{\overset{all}{}}m_j\stackrel{}{a}_j}{\underset{j=1}{\overset{all}{}}m_j}\right)$$
(5)
These new forms of the second law i.e. equations (3) and (5) are invariant under transformation to a more general group of reference frames $`S^{}`$ than Galilean ones which we may call as generalized Galilean transformations.

$$\{\begin{array}{c}t^{}=t\hfill \\ \stackrel{}{a}^{}=\stackrel{}{a}\stackrel{}{b}\hfill \\ \stackrel{}{u}^{}=\stackrel{}{u}\stackrel{}{b}t\stackrel{}{v}\hfill \\ \stackrel{}{x}^{}=\stackrel{}{x}\frac{1}{2}\stackrel{}{b}t^2\stackrel{}{v}t+\stackrel{}{x}_0\hfill \end{array}$$
(6)
Here $`\stackrel{}{b}`$ , $`\stackrel{}{v}`$ and $`\stackrel{}{x}_0`$ are the constant acceleration ,velocity and position of $`S^{}`$ with respect to $`S`$ at $`t=0`$ respectively.

Within the set of reference frames $`S^{}`$ with different values of $`\stackrel{}{b}`$ and $`\stackrel{}{v}`$ there is a unique subset call it $`S^{\prime \prime }`$ moving with acceleration $`\stackrel{}{b}_0`$

$$\stackrel{}{b}_0=\frac{\underset{j=1}{\overset{all}{}}m_j\stackrel{}{a}_j}{\underset{j=1}{\overset{all}{}}m_j}$$
(7)
with respect to $`S`$. In $`S^{\prime \prime }`$ the form of Newton’s second law of inertia is recovered, $`\stackrel{}{F}_i=m\stackrel{}{a}^{\prime \prime }`$. Therefore, inertial frames are a unique set of frames which are moving with acceleration $`\stackrel{}{b}_0`$ with respect to $`S`$ with $`\stackrel{}{v}`$ and $`\stackrel{}{x}_0`$ be any arbitrary value. This again may be accounted for as a feasible formulation of Mach’s idea about inertial frames. Since $`\stackrel{}{b}_0`$ is determined globally then local inhomogeneities have a negligible effect on it and inertial frames.

It is evident that equations (3) and (5) also satisfy Newton’s third law. For a two particle system we have

$$\stackrel{}{F}_1=\stackrel{}{F}_2=\mu .c_1.c_2(\stackrel{}{a}_1\stackrel{}{a}_2)$$
(8)
and for a system with N particles it makes

$$\underset{j=1}{\overset{all}{}}\stackrel{}{F}_i=0$$
(9)
So there is no need to introduce Newton’s third law as an extra principle.

We can extend this model to Newton’s law of gravitation. Equivalence principle here means that the source of inertia and gravitation is the same. Then if we define the gravitational force between two particles of inertial charges $`c_1`$ and $`c_2`$ as

$$\stackrel{}{F}_G=\frac{\mu ^2.c_1.c_2}{\stackrel{}{r}_{12}^2}$$
(10)
where $`\stackrel{}{r}_{12}=\stackrel{}{r}_1\stackrel{}{r}_2`$ is the relative separation between particles 1 and 2. Then we can express gravitational constant $`G`$ in terms of inertial charges $`c_i`$s

$$G=\left(\underset{j=1}{\overset{all}{}}c_j\right)^{\frac{1}{2}}$$
(11)
which may be accounted for as another aspect of Machian implications that the so-called physical constants of nature (like G) are to be determined with some global features of the universe. Since the inertial mass and the gravitational constant are finite quantities, this means that $`\left(\underset{j=1}{\overset{all}{}}c_j\right)`$ is finite too, and the universe cannot be infinitely extended.

The Lagrangean formalism based on this model of inertial law is the same as the one for NM except that the kinematic energy of system T which is equal to $`\underset{i}{}\frac{1}{2}m_iv_{i}^{}{}_{}{}^{2}`$ is replaced by

$`T`$ $`=`$ $`{\displaystyle \underset{i}{}}{\displaystyle \frac{1}{2}}m_iv_{i}^{}{}_{}{}^{2}{\displaystyle \frac{[\underset{i}{}m_i\stackrel{}{v}_i]^2}{2\underset{i}{}m_i}}`$
$`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{i}{}}{\displaystyle \underset{j}{}}m_im_j{\displaystyle \frac{(\stackrel{}{v}_i\stackrel{}{v}_j)^2}{\underset{j}{}m_j}}`$
This is an invariant scalar. Then total energy of the system i.e. the sum of kinetic energy T and potential energy $`V(r_{ij})`$ , a function of relative separation $`r_{ij}=\stackrel{}{r}_i\stackrel{}{r}_j`$, is invariant in all non-rotating reference frames.

Remarks

* \- It may seem that choosing non-rotating frames is some kind of restriction which reduces the generality of the chosen frames. This is however, not the case, since rotating frames with respect to the whole universe are distinguishable and can be fixed.
* \- Our results show that $`G`$ and $`m`$ are not constants of nature but may change anytime the total inertial charge of the universe changes significantly. This can happen in the early stage of evolution of the universe, i.e. in the epochs of pair production.
* \- The concept of energy introduced with this theory is independent of the reference frame.
* \- There is no inertial strucure for empty space.

These may show us how some thing should change in a modified theory of relativity.


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# SOLAR CYCLE INDUCED VARIATIONS IN GONG p-MODE FREQUENCIES AND SPLITTINGS

## 1 INTRODUCTION

It is now well established that the frequency and splitting coefficients of the solar p-mode oscillations vary during the sun’s magnetic cycle. The initial results were obtained from the study of the low degree modes arising from irradiance measurements (Woodard and Noyes, 1985) and observations from the Sun as a star (Elsworth et al., 1990). Libbrecht & Woodard (1990) extended this study for higher $`\mathrm{}`$ modes ( 5 $`<\mathrm{}<`$ 140) and showed that the frequency shift increases with $`\mathrm{}`$ values. On the basis of observations made during the rising phase of cycle 22, Woodard et al. (1991) found a strong correlation between the frequency shift and the total surface magnetic field. The study of solar cycle 22 was extended to the initial phase of cycle 23 by several groups using the Global Oscillation Network Group (GONG) data and MDI data from the Solar and Heliospheric Observatory (SOHO). Anderson, Howe, & Komm (1998) analysed two 3-month equivalent data sets from GONG and MDI during the early phase of the solar cycle 23 and found evidence of small shifts in frequency and $`a_4`$ splitting coefficients. Using the GONG data, Bhatnagar, Jain, & Tripathy (1999) studied the variation in frequencies for a period of two years starting from 1995 August. They also presented evidence that magnetic and radiative activity indices correlate differently with frequency shifts. A similar result was earlier obtained by Bachmann & Brown (1993) for the declining phase of cycle 21 and ascending phase of cycle 22. Howe, Komm, & Hill (1999) have also analysed GONG frequencies over modes near 3 mHz and confirm the solar cycle changes in mean frequencies. In addition to frequency canges, they analysed the splitting coefficients up to $`a_6`$ and reported that the even $`a`$ coefficients show a larger correlation with the corresponding even Legendre component of the magnetic activity than with the average magnetic flux. Centroid frequencies and multiplet frequency splittings obtained from the MDI have also been studied by Dziembowski et al. (1998, 1999) who primarily concentrated on the solar cycle changes of the spherical structure of the Sun.

Recently, Basu & Antia (1999) and Howe, Komm, & Hill (2000) found a small temporal variation in GONG odd-order splitting coefficients upto 15th order. Using different inversion techniques, they found a systematic zonal flow migrating towards lower latitudes during the rising phase of the cycle. This flow can be associated with torsional oscillations as seen in BBSO data from 1986–1990 (Woodard and Libbrecht, 1993a). These zonal variations of the Sun’s differential rotation were first noticed by Kosovichev & Schou (1997) using MDI data. They inferred subsurface flows from the rotational splittings of f-mode frequencies in the range of $`\mathrm{}`$ from 120 to 250 using 144 days time series. Toomre et al. (2000) have confirmed these zonal flows using both p\- and f-mode splittings derived from MDI data.

The changes in frequencies and splittings are believed to be caused by a combination of variations in temperature, magnetic and velocity fields. Goldreich et al. (1991) proposed that the variations in solar eigen-frequencies are related to the perturbations in the surface magnetic fields. However, Kuhn (1998) argues that the frequency changes can not be due to the variation in surface magnetic field alone, because a reasonable explanation of the solar cycle acoustic changes must account for the apparent photometric cycle.

In this paper, we analyse the frequencies and splitting coefficients obtained from the GONG project to extend previous works to all possible frequency and $`\mathrm{}`$ ranges. We also infer the solar rotation rate as a function of depth and latitude using the analytical approach of Morrow (1988) involving the odd order splitting coefficients. Our results reveal a significant temporal variation in the rotation rate at high latitudes over the period of study.

## 2 FREQUENCY SHIFT AND ACTIVITY INDICES

The analysis presented here uses thirty three, 108-days frequency data sets (maximum available at the time of writing) from the GONG project, covering a period of more than 3 years from 1995 May 7 to 1998 October 17. Each data set consists of 3 GONG months (GMs), where each GM extends over 36 days and contains $`m`$–averaged p-mode multiplets. The mode frequencies were estimated using the standard GONG analysis (Hill et al., 1996) and obtained by fitting orthogonal polynomials as defined in Ritzwoller & Lavely (1991):

$$\nu _{n,l,m}=\nu _{n,l}+\underset{s}{}c_{s,n,l}\gamma _{s,l}(m),$$
(1)
where $`\gamma _{s,l}`$($`m`$) are the orthogonal polynomials for given value of degree $`\mathrm{}`$ and $`c_{s,n,l}`$ are the splitting functions. The remaining symbols in Equation (1) have their usual meanings. Thus, this study differs from the earlier analysis of GONG data by Bhatnagar, Jain, & Tripathy (1999) and Howe, Komm, & Hill (1999), wherein the frequencies used were derived from Legendre polynomial expansion.

The mean shift $`\delta \nu `$ for a given $`\mathrm{}`$ and $`n`$ is calculated from the relation

$$\delta \nu (t)=\underset{n,\mathrm{}}{}\frac{\delta \nu _{n,\mathrm{}}(t)}{\sigma _{n,\mathrm{}}^2}/\underset{n,\mathrm{}}{}\frac{1}{\sigma _{n,\mathrm{}}^2},$$
(2)
where $`\delta \nu _{n,\mathrm{}}(t)`$ is the change in the measured frequency for a given $`\mathrm{}`$ and radial order $`n`$ and $`\sigma _{n,\mathrm{}}`$ is the error in the observed frequency. The mean frequency shifts for GM 5–7 (1995 Sep. 23 – 1996 Jan. 13) and GM 33–35 (1998 Jul. 2 – 1998 Oct. 17), which represents low and high activity periods respectively, are given in Table 1 where GM 1–3 is taken as the reference. It is noted that there are very few modes in the frequency range of 1500–1900 $`\mu `$Hz and 3900–4300 $`\mu `$Hz. It is further observed that the errors in $`\delta \nu `$ for these frequency ranges are larger as compared to other ranges.

The variation of $`\nu `$–averaged frequency shift for three different periods scaled with the inverse of the modal mass as a function of $`\mathrm{}`$ is shown in Figure 1, where the frequency shift is calculated from the mean over all relevant data points. It is clear from the panels in Figure 1 that the shifts are constant for $`\mathrm{}`$ between 40 to 80 in the frequency range of 1900–4300 $`\mu `$Hz whereas there is large scatter at both low and high degrees. The temporal evolution of the mean frequency shift for different $`\mathrm{}`$ ranges scaled by the inverse of the modal mass is shown in Figure 2. Both these figures confirm that the frequency shift is a weak function of degree as was pointed out by Libbrecht & Woodard (1990). From GONG data, Howe, Komm, & Hill (1999) had also found no degree dependence for $`\mathrm{}`$ approximately between 50 to 90 around 3mHz.

In Figure 3, we show the variation of $`\mathrm{}`$–averaged frequency shift as a function of frequency for three independent periods with reference to mean over all these periods. It is evident that the frequency shift has a strong dependence on the frequency which increases with increasing frequencies. This is also consistent with the result of Libbrecht & Woodard (1990) and Anguera Gubau et al. (1992). The temporal evolution of frequency shifts over the entire period of three and half years for different frequency ranges is shown in Figure 4. It is noted that the frequency shifts for all frequency ranges are small upto GM 24–26 and then increases rapidly. This behaviour of the shifts can be understood in terms of the rapid increase of solar activity represented in the same figure by means of the 10.7 cm radio flux. It is further noticed that the frequency shifts for all frequency ranges seem to converge at GM 23–25. This particular aspect was investigated in detail using different combination of GMs as reference months. We conclude that this converging feature depends on reference frequencies and has no physical significance.

The mean shift is correlated with different solar activity indices by computing a mean value for each index over the same 108 days interval over which the frequencies were determined. The activity indices considered are: $`R_I,`$ the International sunspot number obtained from the Solar Geophysical Data (SGD); KPMI, Kitt Peak Magnetic Index from Kitt peak full disk magnetograms; SMMF, Stanford Mean Magnetic Field from SGD; MPSI, Magnetic Plage Strength Index from Mount Wilson magnetograms; FI, total flare index from SGD and T. Ataç (1999, private communication); He I, equivalent width of He I 10830 Å line, averaged over the whole disk from Kitt peak; Mg II, core-to-wing ratio of Mg II line at 2800 Å from L. Floyd (1999, private communication); $`F_{10}`$, integrated radio flux at 10.7 cm from SGD. A linear least square fitting is carried out between mean activity indices and frequency shifts. The fitting parameters, $`\chi ^2`$ values and both Pearson’s ($`r_p`$) and Spearman’s ($`r_s`$) correlation coefficients are summarised in Table 2. It is observed that both the magnetic and radiative indices are equally correlated with the frequency shifts in contrast to the earlier results of Bachmann & Brown (1993) and Bhatnagar, Jain, & Tripathy (1999). This aspect needs to be investigated in detail with the availability of more high precision data sets in future. The correlation coefficients for different frequency ranges are given in Table 3. It is observed that the correlation has an increasing trend with frequency ranges and are marginally different at higher frequencies.

## 3 SPLITTING COEFFICIENTS

It is known that the solar differential rotation and other symmetry breaking factors like magnetic field can lift the degeneracy of the solar acoustic modes and split the eigen frequencies. Individual mode splittings can be represented by Legendre polynomial expansion in ($`m/L`$)

$$\nu _{n,\mathrm{},m}=\nu _{n,\mathrm{}}+L\underset{s=1}{\overset{m}{}}a_{s,n,\mathrm{}}P_s(m/L),$$
(3)
where the expansion coefficients $`a_{s,n,\mathrm{}}`$ are known as the splitting coefficients and $`L^2`$ = $`\mathrm{}(\mathrm{}+1)`$. Alternatively, the splitting data can be expressed in terms of the splitting functions, $`c_{s,n,\mathrm{}}`$, as given in Equation (1). These functions are generally converted to a new set of splitting coefficients $`b`$’s by using the relation (Ritzwoller & Lavely, 1991)

$$c_s=\frac{(4\pi )^{1/2}}{(2\mathrm{}+1)L^2H_s^1F_s}b_s,$$
(4)
where

$$H_s^1=(1)^{(\mathrm{}1)}\left(\begin{array}{ccc}s& \mathrm{}& \mathrm{}\\ 0& 1& 1\end{array}\right),$$
(5)
and

$$F_s=\left[\frac{(2\mathrm{}s)!}{(2\mathrm{}+s+1)!}\right]^{1/2}.$$
(6)
The odd-order coefficients measure the solar rotation while the even-order coefficients probe the symmetry about the equator. The nonzero values of these even coefficients reflect the pole–equator asymmetries in the solar structure. It may be noted that we have retained terms up to $`s`$ = 5 because the coefficients and their errors are of similar magnitude for $`s`$ $`>`$ 5 in the data sets used in this study.

The temporal variation of $`b_2`$ and $`b_4`$ coefficients is shown in Figure 5. We find that $`b_2`$ has strong correlation with activity, whereas $`b_4`$ is anticorrelated. The results of linear regression analysis between these coefficients and activity indices are given in Table 4. It may be noted that $`b_2`$ has a better quantitative correlation than $`b_4`$. We observe that during the onset of cycle 23, the change in $`b_4`$ is more prominent than the change in $`b_2`$. Howe, Komm, & Hill (1999) have also studied the temporal variation of the $`b_2`$ and $`b_4`$ coefficients at 3 mHz and reported a phase shift of one year between them. The odd order coefficients, $`b_1`$, $`b_3`$ and $`b_5`$, calculated using the expression (4) are used to obtain information on the variation of the Sun’s rotation rate with depth and latitude.

### 3.1 Solar Rotation Rate

The solar rotation rate using the helioseismic data are generally obtained through two different methods: in forward approach the frequency splittings are computed for a chosen solar rotation model and then compared with the observed splittings. In the inverse method the measured frequency splittings are used directly to produce a single function for the angular velocity. In this study, we use the analytical method of Morrow (1988) where the appropriate combination of odd order splitting coefficients reflects the depth variation of angular velocity at chosen co–latitude $`\varphi `$, (90° $``$ $`\theta `$, where $`\theta `$ is the latitude).

$$\mathrm{\Omega }^{nl}(\varphi )=\underset{i=0}{\overset{2}{}}d_{2i+1}(\varphi )a_{2i+1}^{nl},$$
(7)
where

$`d_1`$ $`=`$ $`1,`$ (8)
$`d_3`$ $`=`$ $`[15cos^2\varphi ],`$ (9)
$`d_5`$ $`=`$ $`[114cos^2\varphi +21cos^4\varphi ].`$ (10)
The corresponding rotation rate at equator is given by

$$\mathrm{\Omega }^{nl}(90^\mathrm{°})=a_1^{nl}+a_3^{nl}+a_5^{nl},$$
(11)
and at pole

$$\mathrm{\Omega }^{nl}(0^\mathrm{°})=a_1^{nl}4a_3^{nl}+8a_5^{nl}.$$
(12)
Wilson, Burtonclay & Li (1997) have recently modified expressions (7-12) by including higher order splitting coefficients up to $`a_7`$. The addition of further terms to these expressions seems to be of limited utility due to the larger relative errors of higher order coefficients. We compute $`a`$’s from $`c`$’s using Equations (61-63) of Ritzwoller & Lavely (1991). In Figure 6, we show the variation of $`a_1`$, $`a_3`$ and $`a_5`$ with degree. The non-zero values of these coefficients indicate that the solar rotation rate changes with depth and latitude.

In Figure 7, we show the sun’s rotation rate at different latitudes as a function of depth approximated by $`\nu /L`$ which corresponds to the radius of lower turning point $`r_t`$ as defined by the relation (Christensen-Dalsgaard & Berthomieu, 1991)

$$r_t=\frac{c(r_t)}{2\pi }\frac{L}{\nu }$$
(13)
where $`c`$ is the sound speed. A higher value of $`\nu /L`$ denotes a smaller value of $`r_t`$ and hence a greater depth. As a result the low degree modes are sensitive to the rotation from surface to the core, while the high degree modes probe solar rotation near the surface. From Figure 7, we find that for $`\theta `$ = 0° and 30°, the rotation rate first increases and then decreases as depth increases. On the other hand, the rotation rate for $`\theta `$ = 45° and 60° remains constant below the base of the convection zone ( $`\nu /L`$ $``$ 70), beyond which it starts increasing. We further notice that for $`\theta `$ = $`60^\mathrm{°}`$, the increase in internal rotation rate above $`\nu /L`$ = 70, is steeper than at $`\theta `$ = $`45^\mathrm{°}`$. This suggests that, as the depth increases, the equatorial rotation rate decreases while the polar rotation rate increases. Wilson, Burtonclay & Li (1997) have analysed frequency splittings derived from first four months of GONG data and showed that the greatest changes in rotation occurs in the region below the base of the convection zone. They also found that the data does not support models exhibiting a discontinuous shear between the convection zone and a uniformly rotating radiative region. Thompson et al. (1996) have also estimated the solar rotation rate using the first few month’s of GONG data. Our results derived from analytical expressions are in gross agreement with earlier studies.

The observed rotation rate is conventionally expressed in terms of even powers of cos$`\varphi `$:

$$\mathrm{\Omega }(r,\varphi )=A+Bcos^2\varphi +Ccos^4\varphi ,$$
(14)
where $`\varphi `$ as defined earlier is the solar co–latitude. The surface constants $`A`$, $`B`$, and $`C`$ are related to splitting coefficients by the relations

$`a_1`$ $`=`$ $`A+{\displaystyle \frac{1}{5}}B+{\displaystyle \frac{3}{35}}C,`$ (15)
$`a_3`$ $`=`$ $`{\displaystyle \frac{1}{5}}B{\displaystyle \frac{2}{15}}C,`$ (16)
$`a_5`$ $`=`$ $`{\displaystyle \frac{1}{21}}C.`$ (17)
Various workers have calculated these surface coefficients using different data sets. Brown et al. (1989) found $`A`$ = 462.8 nHz, $`B`$ = $``$ 56.7 nHz and $`C`$$``$ 75.9 nHz for $`r`$0.723 $`R_{}`$ using CaII K intensity data taken from South pole. Based on 100-day observations made at BBSO for $`\mathrm{}`$ between 10 to 60, Libbrecht (1989) found the best fit with $`A`$ = 461 nHz, $`B`$ = $``$ 60.5 nHz and $`C`$ = $``$ 75.4 nHz. We have derived these constants from the GONG data for the 20 $``$ $`\mathrm{}`$ $``$ 100 and obtained the best fit with $`A`$= 459.54 nHz, $`B`$ = $``$ 61.90 nHz and $`C`$ = $``$ 70.20 nHz. Using the odd-order splitting coefficients of GONG data derived from Legendre polynomial expansion, we obtain, $`A`$ = 459.97 nHz, $`B`$ = $``$ 61.33 nHz, and $`C`$ = $``$ 70.93 nHz. The values obtained for $`A`$, $`B`$ and $`C`$ using GONG data are in close agreement with the earlier values.

In Figure 8, we plot the average surface rotation rate as a function of latitude using the derived coefficients. In the same plot, we have also shown the results from inversion techniques (Antia, private communication) and Doppler surface measurements (Snodgrass, 1984). It is clear that the rotation rate changes significantly in mid latitude while the change in rotation rate near pole and equator is small. We find that the rotation rate derived from the GONG data agrees well with other results. However, the inverted rotation profile departs from other rotation rates beyond the latitude of 70°, probably due to the resolution limitation in inversion techniques.

Woodard and Libbrecht (1993b) studied the time variation of equatorial rotation rate for different $`\mathrm{}`$ ranges and suggested a small variation from year to year. To investigate temporal variations, Figure 9 shows the residual surface rate at four different latitudes, obtained after subtracting the average angular velocity shown in Figure 8 by long dash line. It is evident that at higher latitudes, the residual rotation rate, commonly known as zonal flows, is time dependent. This has a magnitude of approximately 2 nHz at a latitude of 60°. Recently Basu & Antia (1999), Howe, Komm, & Hill (2000) using GONG data, and Toomre et al. (2000) using MDI data, have also reported a small but significant time variation in the rotation rate. Our results derived from the analytical method are consistent with these inversion studies.

## 4 CONCLUSION

The mean frequency shift varies over the solar cycle and is correlated similarly with both the magnetic and radiative indices.

The $`b_2`$ splitting coefficient is linearly correlated with the activity indices, while $`b_4`$ is anti-correlated for all degree and frequency ranges.

We detect a small but significant variation in the rotation rate derived from the linear combination of odd order coefficients over a period of three and half years.

We thank the anonymous referee for his critical remarks which impreved the manuscript. We also thank T. Ataç, L. Floyd, and R.K. Ulrich for supplying us the Flare index, Mg II, and MPSI data respectively and H. M. Antia for providing us the inversion results. This work utilises data obtained by the Global Oscillation Network Group project, managed by the National Solar Observatory, a Division of the National Optical Astronomy Observatories, which is operated by AURA, Inc. under cooperative agreement with the National Science Foundation. The data were acquired by instruments operated by Big Bear Solar Observatory, High Altitude Observatory, Learmonth Solar Obsrvatory, Udaipur Solar Observatory, Instituto de Astrophsico de Canaris, and Cerro Tololo Interamerican Observatory. NSO/Kitt Peak magnetic, and Helium measurements used here are produced cooperatively by NSF/NOAO; NASA/GSFC and NOAA/SEL. This work is partially supported under the CSIR Emeritus Scientist Scheme and Indo–US collaborative programme–NSF Grant INT–9710279.


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# Wave packets and initial conditions in quantum cosmology

## 1 Introduction

The question of the construction and interpretation of wave packets in quantum cosmology has been attracting much attention in recent years. Numerous studies have been done to obtain a quantum theory for gravity and to understand its connection with classical physics. This interest is due to our desire for obtaining the basic dynamical laws for predicting the large scale structure of the universe and the determination of the initial conditions from which the universe has evolved. The latter is of particular importance since if we accept that the fundamental laws of physics are quantum mechanical in nature, then the question of initial conditions is embodied in finding the initial condition for the quantum cosmological state of the universe, see for example . The search for the fundamental dynamical laws has been underway since Newtonian times but that for the initial conditions is relatively recent and great efforts have been made in this direction. This is, of course, not surprising as the underlying theories for describing the fundamental laws of physics are local quantum field theories above the Plank length. However, theories explaining the initial quantum state of the universe are non-local. They imply regularities in space that have emerged on large cosmological scales. It is only in recent years that the progress in observational cosmology on large enough scales of both space and time may present us with the opportunity to look at the viability of the predictions of a theory of the initial state of the universe.

The problem of the relation between quantum cosmology and classical physics is an important one that exists even in simple models . Most authors consider semi-classical approximations to the Wheeler-DeWitt (WD) equation and refer to regions in configuration space where these solutions are oscillatory or exponentially decaying as representing classically allowed or forbidden regions, respectively. These regions are mainly determined by the initial conditions for the wave function. Two popular proposals for the initial conditions are the no boundary proposal  and the spontaneous nucleation from nothing . These proposals have been attractive to many authors because they lead to some classes of classical solutions represented by certain trajectories which posses important features such as predicting an inflationary phase.

In quantum cosmology, in analogy with ordinary quantum mechanics, one is generally concerned with the construction of wave functions by the superposition of the ‘energy eigenstates’ which would peak around the classical trajectories in configuration space. However, contrary to ordinary quantum mechanics, a parameter describing time is absent in quantum cosmology so that the initial conditions would have to be expressed in terms of an intrinsic time parameter, which in the case of the WD equation could be taken as the local scale factor for the three geometry . Also, since the sign of the kinetic term for the scale factor is negative, a formulation of the Cauchy problem for the WD equation is possible. The existence of such a sign is one of the exclusive features of gravity with many other interesting implications.

The construction of wave packets resulting from the solutions of the WD equation has been a feature common to most of the recent research work in quantum cosmology . In particular, in reference  the construction of wave packets in a Friedmann universe is presented in detail and appropriate boundary conditions are motivated. Generally speaking one of the aims of these investigations has been to find wave packets whose probability distributions coincide with the classical paths obtained in classical cosmology. In these works, the authors usually consider model theories in which a self interacting scalar field is coupled to gravity in a Robertson-Walker type universe. The resulting WD equation is often in the form of an isotropic oscillator-ghost-oscillator and is separable in the configuration space variables. The general solution can thus be written as a sum over the product of simple harmonic oscillator wave functions with different frequencies. As usual, the coefficients in the sum are chosen according to the initial conditions, which are usually specified by giving the wave function at $`R=0`$, where $`R`$ is the scale factor. The choice of the initial conditions, or equivalently the coefficients in the sum, varies considerably amongst various authors. In , the wave packets are constructed by summing over the first few terms with the coefficients being chosen arbitrarily to be equal. However, in  the initial state is chosen to be a Gaussian with the appropriate symmetry and the undetermined coefficients being set to zero. The resulting wave packets, in general, lack one or more of the following desired properties. Firstly, one expects the wave packet describing a physical system to posses a certain degree of smoothness. Secondly, there should be a good classical-quantum correspondence, which means that not only the wave packet should be centered around the classical path, but also the crest of the probability distribution should coincide as closely as possible with the classical trajectory. Also, to each distinct classical path there should correspond a unique wave packet.

The purpose of this paper is to address and clarify some of these issues. We first argue that specifying the wave function at $`R=0`$ does not furnish a complete set of initial conditions. We then suggest an ansatz for the expansion coefficients, or equivalently for the initial conditions, which produces wave packets with the desired properties mentioned above. The organization of the paper is a follows: In section two we give a short review of some of the models that lead to the class of WD equations considered here. Section three deals with the construction of the wave packets and in section four we discuss the salient features of this work and conclusions are drawn.

## 2 Robertson-Walker cosmology with coupled scalar field

In Robertson-Walker cosmology, one often considers the coupling of a scalar field to gravity. The resulting field equations include a ‘zero energy constraint’. The WD equation in quantum cosmology is the result of the quantization consistent with this constraint. In some simple minisuperspace models the solutions to the WD equation describe a system of oscillators with unequal or equal frequencies. It would therefore be appropriate at this point to give a brief review of some of the minisuperspace models leading to the class of WD equations whose solutions are to be discussed here.

Consider Einstein’s field equation coupled to a scalar field

$`R_{\mu \nu }{\displaystyle \frac{1}{2}}g_{\mu \nu }R`$ $`=`$ $`\kappa T_{\mu \nu }(\varphi ),`$ (1)
$`\mathrm{\Delta }^2\varphi {\displaystyle \frac{U}{\varphi }}`$ $`=`$ $`0,`$ (2)
where $`T_{\mu \nu }`$ is the energy-momentum tensor coupling the scalar field $`\varphi `$ to gravity and $`U`$ is a scalar potential for the scalar field $`\varphi `$. Hereafter, we shall work in a system of units in which $`\kappa =1`$. We parameterize the metric as

$`g=dt^2+R^2(t){\displaystyle \frac{dr^idr^i}{(1+kr^2/4)^2}},`$ (3)
where $`R(t)`$ is the usual scale factor and $`k=+1,0,1`$ corresponds to a closed, flat or open universe respectively. The field equations resulting from (1) and (2) with the metric given by (3) can be written as

$`3\left[\left({\displaystyle \frac{\dot{R}}{R}}\right)^2+{\displaystyle \frac{k}{R^2}}\right]`$ $`=`$ $`{\displaystyle \frac{\dot{\varphi }^2}{2}}+U(\varphi ),`$ (4)
$`2\left({\displaystyle \frac{\ddot{R}}{R}}\right)+\left({\displaystyle \frac{\dot{R}}{R}}\right)^2+{\displaystyle \frac{k}{R^2}}`$ $`=`$ $`{\displaystyle \frac{\dot{\varphi }^2}{2}}+U(\varphi ),`$ (5)
$`\ddot{\varphi }+3{\displaystyle \frac{\dot{R}}{R}}\dot{\varphi }+{\displaystyle \frac{U}{\varphi }}`$ $`=`$ $`0.`$ (6)
Here a dot represents differentiation with respect to time. The choice of $`U(\varphi )`$ is important and a particularly interesting choice with three free parameters is

$`U(\varphi )=\lambda +{\displaystyle \frac{m^2}{2\alpha ^2}}\mathrm{sinh}^2(\alpha \varphi )+{\displaystyle \frac{b}{2\alpha ^2}}\mathrm{sinh}(2\alpha \varphi ).`$ (7)
In the above expression $`\lambda `$ may be identified with a cosmological constant, $`m^2=^2U/\varphi ^2|_{\varphi =0}`$ is a mass squared parameter, $`b`$ is a coupling constant and $`\alpha ^2=\frac{3}{8}`$. The Lagrangian giving the above equations of motion can be written as

$`L=3R\dot{R}^2+3kR+R^3[\dot{\varphi }^2/2U(\varphi )].`$ (8)
This Lagrangian can be cast into a simple form for $`k=0`$ describing an oscillator-ghost- oscillator system . Briefly, this is achieved by using the transformations $`X=R^{3/2}\mathrm{cosh}(\alpha \varphi )`$ and $`Y=R^{3/2}\mathrm{sinh}(\alpha \varphi )`$, which transform the term $`R^3U(\varphi )`$ into a quadratic form. Upon using a second transformation to eliminate the coupling term in the quadratic form, we arrive at new variables $`u`$ and $`v`$, which are linear combinations of $`X`$ and $`Y`$, in terms of which the Lagrangian takes on the simple form

$`L(u,v)={\displaystyle \frac{1}{2}}\left[(\dot{u}^2\omega _1^2u^2)(\dot{v}^2\omega _2^2v^2)\right],`$ (9)
where $`\omega _{1,2}^2=3\lambda /4+m^2/2\sqrt{m^44b^2}/2`$. The classical solutions for $`k0`$ are given in . The corresponding quantum cosmology for $`k=0`$ is described by the WD equation written as

$`H\psi (u,v)=\left\{{\displaystyle \frac{^2}{u^2}}+{\displaystyle \frac{^2}{v^2}}+\omega _1^2u^2\omega _2^2v^2\right\}\psi (u,v)=0.`$ (10)
Alternatively, one may obtain a minisuperspace model described by the same WD equation but with equal frequencies by considering a closed Robertson-Walker universe with a vanishing cosmological constant and containing a conformally coupled scalar field . A different model leading to the same WD equation with equal frequencies is obtained by considering a Kaluza-Klein cosmology with a negative cosmological constant described by the metric

$`g=dt^2+R^2(t){\displaystyle \frac{dr^idr^i}{(1+kr^2/4)^2}}+a^2(t)d\rho ^2,`$ (11)
where $`a(t)`$ is the radius of the compactified space . The Lagrangian describing the above cosmology is

$`L={\displaystyle \frac{1}{2}}Ra\dot{R}^2{\displaystyle \frac{1}{2}}R^2\dot{R}\dot{a}+{\displaystyle \frac{1}{2}}kRa{\displaystyle \frac{1}{6}}\mathrm{\Lambda }R^3a,`$ (12)
where $`\mathrm{\Lambda }`$ is the cosmological constant. By defining $`\omega ^2=\frac{2}{3}\mathrm{\Lambda }`$ and changing the variables as

$`u={\displaystyle \frac{1}{\sqrt{8}}}\left[R^2Ra{\displaystyle \frac{3k}{\mathrm{\Lambda }}}\right],v={\displaystyle \frac{1}{\sqrt{8}}}\left[R^2+Ra{\displaystyle \frac{3k}{\mathrm{\Lambda }}}\right],`$ (13)
$`L`$ takes on the form

$`L={\displaystyle \frac{1}{2}}\left[(\dot{u}^2\omega ^2u^2)(\dot{v}^2\omega ^2v^2)\right].`$ (14)
It is easily seen that the corresponding quantum cosmology is the same as (10) with equal frequencies. We therefore focus attention on solutions to equation (10).

## 3 Wave packets

Equation (10) is separable in the minisuperspace variables and a solution can be written as

$`\psi _{n_1,n_2}(u,v)=\alpha _{n_1}(u)\beta _{n_2}(v),`$ (15)
where

$`\alpha _n(u)=\left({\displaystyle \frac{\omega _1}{\pi }}\right)^{1/4}{\displaystyle \frac{H_n(\sqrt{\omega _1}u)}{\sqrt{2^nn!}}}e^{\omega _1u^2/2},`$ (16)
$`\beta _n(v)=\left({\displaystyle \frac{\omega _2}{\pi }}\right)^{1/4}{\displaystyle \frac{H_n(\sqrt{\omega _2}v)}{\sqrt{2^nn!}}}e^{\omega _2v^2/2}.`$ (17)
Here $`H_n(x)`$ are Hermite polynomials, and the zero energy condition leads to

$`\left(n_1+{\displaystyle \frac{1}{2}}\right)\omega _1=\left(n_2+{\displaystyle \frac{1}{2}}\right)\omega _2,n_1,n_2=0,1,2,\mathrm{}.`$ (18)
The set $`\left\{\psi _{n_1,n_2}(u,v)\right\}`$ span the zero sector subspace of the Hilbert space of $`^2`$ of measurable square integrable functions on $`𝐑^\mathrm{𝟐}`$ with an inner product defined in the usual way giving

$$\psi _{n_1,n_2}(u,v)\psi _{n_1^{},n_2^{}}(u,v)𝑑u𝑑v=\delta _{n_1,n_1^{}}\delta _{n_2,n_2^{}}.$$

That is, the orthonormality and completeness of the basis functions follow from those of the Hermite polynomials.

We can construct a general wave packet as follows

$`\psi (u,v)={\displaystyle \underset{n_1,n_2}{}^{}}A_{n_1,n_2}\alpha _{n_1}(u)\beta _{n_2}(v),`$ (19)
where the prime on the sum indicates summing over all values of $`n_1`$ and $`n_2`$ satisfying the constraint (18). As the signs of the kinetic terms in equation (10) indicate, we can take $`v`$ as playing the role of the scale factor and hence the initial condition on $`\psi `$ is specified by

$$\psi (u,0)=\underset{n_1}{}^{\prime \prime }c_{n_1}\alpha _{n_1}(u),$$
(20)
where the coefficients $`c_{n_1}`$ are arbitrary and the double prime on the sum indicates that the sum runs only over those values of $`n_1`$ satisfying the constraint (18) with $`n_2`$ being even. The constraint on the evenness of $`n_2`$ is due to the fact that $`\beta _{n_2}(0)=0`$ for odd values of $`n_2`$. We choose the coefficients $`c_n`$ to be the same as those of the coherent states of a one dimensional simple harmonic oscillator, that is

$$c_n=e^{\frac{1}{4}|\chi _0|^2}\frac{\chi _0^n}{\sqrt{2^nn!}},$$
(21)
where $`\chi _0`$ is an arbitrary complex number. Comparing the coefficients in equations (19,20) one finds for even values of $`n_2`$

$`\left({\displaystyle \frac{\omega _2}{\pi }}\right)^{1/4}{\displaystyle \frac{A_{n_1,n_2}}{\sqrt{2^{n_2}n_2!}}}`$ $`=`$ $`{\displaystyle \frac{c_{n_1}}{H_{n_2}(0)}}`$ (22)
$`=`$ $`{\displaystyle \frac{(n_2/2)!c_{n_1}}{(1)^{n_2/2}n_2!}},`$ (23)
and $`A_{n_1,n_2}`$ are arbitrary for odd values of $`n_2`$. This arbitrariness is a consequence of not having specified $`\psi (u,v)/v|_{v=0}`$. The WD equation is a second order PDE, and to have a unique solution one has to specify both the wave function and its derivative at a given point. In the first reference in  the author has chosen $`A_{n_1,n_2}=0`$ for odd values of $`n_2`$, which is equivalent to choosing $`\psi (u,v)/v|_{v=0}=0`$. This might seem to be the only choice when considering equation (22). Here we argue that this is not the best or the ‘canonical’ choice. We have chosen $`A_{n_1,n_2}`$ to be nonzero and to have the same functional form for both even and odd values of $`n_2`$. This can be easily done using equation (23).

The classical paths corresponding to these solutions are the generalized Lissajous ellipsis which have the following parametric representation

$`u(t)=u_0\mathrm{cos}(\omega _1t\theta _0),v(t)=v_0\mathrm{sin}(\omega _2t),`$ (24)
where the zero energy condition demands $`\omega _1u_0=\omega _2v_0`$, and $`\theta _0`$ is an arbitrary phase factor. The classical-quantum correspondence is established by the following equation

$`\chi _0=\sqrt{{\displaystyle \frac{\omega _1}{\omega _2}}}u_0e^{i\theta _0}.`$ (25)
Below we will consider several illustrative examples in order to compare our results with some of the previous works. Let us first consider the simplest case which is when $`\omega _1=\omega _2=\omega `$. Upon inspection of equations (16) through (19), we immediately recognize that, since $`n_1=n_2=n`$ and $`H_n(0)=0`$ for odd values of $`n`$, firstly $`A_{n,n}`$ are not determined for odd values of $`n`$, and secondly

$`\psi (u,0)`$ $`=`$ $`\psi (u,0),`$ (26)
$`{\displaystyle \frac{}{v}}\psi (u,v)|_{v=0}`$ $`=`$ $`{\displaystyle \frac{}{v}}\psi (u,v)|_{v=0}.`$ (27)
Therefore the initial state has to be symmetric and we choose it to be two symmetric Gaussians. However this is automatically taken into account by the restrictions imposed on the sum in equation (20) and the choice of coefficients given in equation (21). Also note that equation (27) is redundant due to the properties of the Hermite polynomials.

Figure 1 Shows the square of the wave packet $`|\psi (u,v)|^2`$ for $`|\chi _0|=12`$ and $`\theta _0=0`$, summing over both even and odd values of $`n`$, and the contour plot of the same wave packet along with the classical path superimposed on it. In practice, in order to obtain a reasonable graphical representation, the minimum number of terms included in the sum, denoted by $`n_{max}`$, has to be of the order of $`|\chi _0|^2`$. As is apparent from these figures, the crest of the wave packet follows the classical path exactly. This is to be compared with figure 1 of the first reference in  which we reproduce here as figure 2 for the same value of $`n_{max}`$ as in figure 1.

At this stage, due to the symmetry of the circle, the only advantage of our choice of the coefficients seems to be that the resulting wave packet is smoother. There is another advantage which will show up in the asymmetric case considered below.

One can easily interpolate between a circle and a line segment simply by changing the value of $`\theta _0`$ between zero and $`\pm \pi /2`$. Figure 3 shows the wave packet for $`|\chi _0|=12`$ and $`\theta _0=1`$.

This is to be compared with figure 3 of the first reference in , where both wave packets with $`\theta _0=\pm 1`$ are superimposed due to the choice $`A_{n,n}=0`$ for odd values of $`n`$. This is not a desired result since in order to establish a classical-quantum correspondence we need to have a separate wave packet for each distinct classical path. At the classical level, a given set of initial conditions, e.g. specifying $`u(0),v(0)`$, uniquely determines a classical path. At the quantum level the initial conditions are specified for example by giving $`\psi (u,0)`$ and $`\psi (u,v)/v|_{v=0}`$, and this will uniquely determine a wave packet. Of course the initial conditions have to be consistent with the symmetries derived from the general structure of the solutions given by equation (19). In the present context, the first condition determines $`A_{n,n}`$ for even values of $`n`$, and the second the odd coefficients. Therefore the simultaneous appearance of two wave packets corresponding to two different classical paths is solely due to choosing $`\psi (u,v)/v|_{v=0}=0`$ or equivalently setting $`A_{n,n}=0`$ for odd values of $`n`$. This is in contrast to the interpretation given by the author in  where he relates the simultaneous appearance of two distinct branches to the superposition principle, and their only means of separation to the decoherence mechanism. As we have shown, for the class of problems discussed here, the use of the ansatz for the initial conditions given above would result in the selection of a unique branch without resorting to the decoherence mechanism.

As $`\theta _0\pm \pi /2`$, the semi-minor axis of the ellipse goes to zero. The results are shown in figure 4. This configuration is to be compared with the approximate one obtained by adding a few terms with equal coefficients in .

Now we discuss the case with unequal frequency and to be concrete we give an example with $`\omega _1=3\omega _2`$. Then the constraint equation (18) gives $`n_2=3n_1+1`$. Therefore at $`v=0`$ only the odd values of $`n_1`$ contribute. Hence we have

$`\psi (u,0)`$ $`=`$ $`\psi (u,0),`$ (28)
$`{\displaystyle \frac{}{v}}\psi (u,v)|_{v=0}`$ $`=`$ $`{\displaystyle \frac{}{v}}\psi (u,v)|_{v=0}.`$ (29)
Again equation (29) is redundant and equation (28) is automatically taken into account by the restrictions imposed on the sum in equation (20). This together with the choice of the coefficients given in equation (21) gives two antisymmetric Gaussians. Only the coefficients $`A_{n_1,n_2}`$ for odd values of $`n_1`$ are determined and those for even values of $`n_1`$ are fixed by the prescription given above. Figure 5 shows a plot of $`|\psi (u,v)|^2`$ for $`|\chi _0|=6`$ and $`\theta _0=0`$. There are pronounced oscillations where there is an overlap between different branches of the wave packet due to the interference effect. This can be clearly seen from a parametric plot of $`|\psi (u,v)|^2`$ along the classical path similar to figure 2, which we do not show here. In figure 5 these oscillations are smoothed out and show up as the bump in the center of the figure. Also note that there is a good classical-quantum correspondence.

Let us next consider the case $`\theta _0=\pi /2`$ whose wave packet is shown in figure 6.

Note that there are pervasive pronounced oscillations due to the overlap. This configuration is to be compared with the approximate one obtained by adding a few terms with equal coefficients in . Moreover, the classical path presented in equation (33) of  , describes an open curve which seems not to be correct. The correct parametric representation of the classical path which is a closed curve is given by equation (24).

## 4 Conclusions

We have described various classical cosmological models leading to a class of WD equations represented by equation (10). The general solution to this equation is a superposition of products of oscillator wave functions, as given in equation (19). We have argued that there exists a set of coefficients which produce wave packets with the following two desired properties. First, they posses a certain degree of smoothness, except when there is an overlap between different branches of the wave packet, where we expect some interference causing genuine pronounced oscillations. Second, there is acceptable classical-quantum correspondence. This correspondence consists of three conditions: The wave packet should have compact support centered around the classical path, the crest of the wave packet should follow as closely as possible the classical path, and to each distinct classical path there should correspond a unique wave packet with the above properties. The set of coefficients presented here as an ansatz is built from the coefficients of the coherent states as in ordinary quantum mechanics and coefficients derived from the functional form of $`H_n(0)`$ for even values of $`n`$. The resulting wave packets can be considered as coherent states in quantum cosmology and they posses all the aforementioned desired properties, and the classical paths are general Lissajous figures. We have shown these features explicitly in several examples. It is worth mentioning that, for the class of problems discussed here, the accomplishment of the last condition for classical-quantum correspondence by our choice of initial conditions shows that there is no need to employ decoherence mechanism to separate out wave packets corresponding to distinct classical paths. However, it is generally accepted that decoherence mechanism plays a major role in quantum-to-classical transition, see  for reviews. Although we have not shown the results explicitly here, we can report that as $`|\chi _0|`$ decreases the degree of classical-quantum correspondence diminishes. Also note that there is a difference between these coherent states and those of quantum mechanics in that the former are static.

We have also shown that there exists a continuous parameter $`\pi /2<\theta _0<\pi /2`$, whose variation gives a continuous interpolation between wave packets or classical paths corresponding to different initial conditions. When $`\theta _0=0`$ we have the most symmetric configurations, i.e. loops which are symmetric about the line $`u=0`$. When $`\theta _0=\pm \pi /2`$ we have the most asymmetric configuration, i.e. curve segments.

This work is also related to the topic of initial conditions in the context of signature transition. For example, at the classical level there is a controversy on the value of $`\dot{R}`$ and $`\dot{\varphi }`$ at the hypersurface of signature transition. The authors of reference  argue that these quantities need to be only continuous across the hypersurface, while in  it is insisted that in addition to the requirement of continuity, the values of both these quantities should be zero at the hypersurface. This choice corresponds to the most asymmetric case $`\theta _0=\pm \pi /2`$, which gives curve segments as classical paths. We speculate that these are unlikely configurations due to their pronounced pervasive oscillations for $`\omega _1/\omega _2>1`$.

AcknowledgementThe authors would like to thank the office of research of Shahid Beheshti University for financial support.


text:
# One-Dimensional Peg Solitaire

## 1 Introduction

Peg solitaire is a game for one player. Each move consists of hopping a peg over another one, which is removed. The goal is to reduce the board to a single peg. The best-known forms of the game take place on cross-shaped or triangular boards, and it has been marketed as “Puzzle Pegs” and “Hi-Q.” Discussions and various solutions can be found in .

In , Guy proposes one-dimensional peg solitaire as an open problem in the field of combinatorial games. Here we show that the set of solvable configurations forms a regular language, i.e. it can be recognized by a finite-state automaton. In fact, this was already shown in 1991 by Plambeck (, Introduction and Ch.5) and appeared as an exercise in a 1974 book of Manna . More generally, B. Ravikumar showed that the set of solvable configurations on rectangular boards of any finite width is regular , although finding an explicit grammar seems to be difficult on boards of width greater than 2.

Thus there is little new about this result. However, it seems not to have appeared in print, so here it is.

###### Theorem 1.1

The set of configurations that can be reduced to a single peg is the regular language $`0^{}L0^{}`$ where

$`L`$ $`=`$ $`1+011+110`$ (1)
$`+`$ $`11(01)^{}\left[00+00(11)^++(11)^+00+(11)^{}1011+1101(11)^{}\right](10)^{}11`$
$`+`$ $`11(01)^{}(11)^{}01+10(11)^{}(10)^{}11.`$
Here $`1`$ and $`0`$ indicate a peg and a hole respectively, $`w^{}`$ means ‘0 or more repetitions of $`w`$’, and $`w^+=ww^{}`$ means ‘1 or more repetitions of $`w`$.’

###### Proof

To prove the theorem, we follow Leibnitz  in starting with a single peg, which we denote

$$1$$

and playing the game in reverse. The first ‘unhop’ produces

$$011\text{ or }110$$

and the next

$$1101\text{ or }1011.$$

(As it turns out, $`11`$ is the only configuration that cannot be reduced to a single peg without using a hole outside the initial set of pegs. Therefore, for all larger configurations we can ignore the $`0`$’s on each end.)

We take the second of these as our example. It has two ends, $`10\mathrm{}`$ and $`\mathrm{}11`$. The latter can propagate itself indefinitely by unhopping to the right,

$$1010101011.$$

When the former unhops, two things happen; it becomes an end of the form $`11\mathrm{}`$ and it leaves behind a space of two adjacent holes,

$$110010101011.$$

Furthermore, this is the only way to create a $`00`$. We can move the $`00`$ to the right by unhopping pegs into it,

$$111111110011.$$

However, since this leaves a solid block of $`1`$’s to its left, we cannot move the $`00`$ back to the left. Any attempt to do so reduces it to a single hole,

$$111111101111.$$

Here we are using the fact that if a peg has another peg to its left, it can never unhop to its left. We prove this by induction: assume it is true for pairs of pegs farther left in the configuration. Since adding a peg never helps another peg unhop, we can assume that the two pegs have nothing but holes to their left. Unhopping the leftmost peg then produces $`1101`$, and the original (rightmost) peg is still blocked, this time by a peg which itself cannot move for the same reason.

In fact, there can never be more than one $`00`$, and there is no need to create one more than once, since after creating the first one the only way to create another end of the form $`10\mathrm{}`$ or $`\mathrm{}01`$ is to move the $`00`$ all the way through to the other side

$$111111111101$$

and another $`00`$ created on the right end now might as well be the same one.

We can summarize, and say that any configuration with three or more pegs that can be reduced to a single peg can be obtained in reverse from a single peg by going through the following stages, or their mirror reflection:

1. We start with $`1011`$. By unhopping the rightmost peg, we obtain $`10(10)^{}11`$. If we like, we then
2. Unhop the leftmost peg one or more times, creating a pair of holes and obtaining $`11(01)^{}00(10)^{}11`$. We can then
3. Move the $`00`$ to the right (say), obtaining $`11(01)^{}(11)^{}00(10)^{}11`$. We can stop here, or
4. Move the $`00`$ all the way to the right, obtaining $`11(01)^{}(11)^{}01`$, or
5. Fill the pair by unhopping from the left, obtaining $`11(01)^{}(11)^{}1011(10)^{}11`$.

Equation 1 simply states that the set of configurations is the union of all of these plus $`1`$, $`011`$, and $`110`$, with as many additional holes on either side as we like. Then $`0^{}L0^{}`$ is regular since it can be described by a regular expression , i.e. a finite expression using the operators $`+`$ and $``$. ∎

Among other things, Theorem 1.1 allows us to calculate the number of distinct configurations with $`n`$ pegs, which is

$$N(n)=\{\begin{array}{ccc}1\hfill & & n=1\hfill \\ 1\hfill & & n=2\hfill \\ 2\hfill & & n=3\hfill \\ 157n+n^2\hfill & & n4,n\text{ even}\hfill \\ 167n+n^2\hfill & & n5,n\text{ odd}\hfill \end{array}$$

Here we decline to count $`011`$ and $`110`$ as separate configurations, since many configurations have more than one way to reduce them.

We also have the corollary

###### Corollary 1

There is a linear-time strategy for playing peg solitaire in one dimension.

###### Proof

Our proof of Theorem 1.1 is constructive in that it tells us how to unhop from a single peg to any feasible configuration. We simply reverse this series of moves to play the game. ∎

More generally, a configuration that can be reduced to $`k`$ pegs must belong to the regular language $`(0^{}L0^{})^k`$, since unhopping cannot interleave the pegs coming from different origins . This leads to the following algorithm:

###### Theorem 1.2

There is a linear-time strategy for reducing any one-dimensional peg solitaire configuration to the minimum possible number of pegs.

###### Proof

Suppose we are given a string $`c_0c_1c_2\mathrm{}c_{n1}`$ where each $`c_i\{0,1\}`$. Let $`𝒜`$ be a nondeterministic finite automaton (without $`ϵ`$-transitions) for $`0^{}L0^{}`$, where $`A`$ is the set of states in $`𝒜`$, $`s`$ is the start state, and $`T`$ is the set of accepting states. We then construct a directed acyclic graph $`G`$ as follows: Let the vertices of $`G`$ consist of all pairs $`(a,i)`$ where $`aA`$ and $`0in`$. Draw an arc from $`(a,i)`$ to $`(b,i+1)`$ in $`G`$ whenever $`𝒜`$ makes a transition from state $`a`$ to state $`b`$ on symbol $`c_i`$. Also, draw an arc from $`(t,i)`$ to $`(s,i)`$ for any $`tT`$ and any $`0in`$. Since $`|𝒜|=𝒪(1)`$, $`|G|=𝒪(n)`$.

Then any path from $`(s,0)`$ to $`(s,n)`$ in $`G`$ consists of $`n`$ arcs of the form $`(a,i)`$ to $`(b,i+1)`$, together with some number $`k`$ of arcs of the form $`(t,i)`$ to $`(s,i)`$. Breaking the path into subpaths by removing all but the last arc of this second type corresponds to partitioning the input string into substrings of the form $`0^{}L0^{}`$, so the length of the shortest path from $`(s,0)`$ to $`(s,n)`$ in $`G`$ is $`n+k`$, where $`k`$ is the minimum number of pegs to which the initial configuration can be reduced. Since $`G`$ is a directed acyclic graph, we can find shortest paths from $`(s,0)`$ by scanning the vertices $`(a,i)`$ in order by $`i`$, resolving ties among vertices with equal $`i`$ by scanning vertices $`(t,i)`$ (with $`tT`$) earlier than vertex $`(s,i)`$. When we scan a vertex, we compute its distance to $`(s,0)`$ as one plus the minimum distance of any predecessor of the vertex. If the vertex is $`(s,0)`$ itself, the distance is zero, and all other vertices $`(a,0)`$ have no predecessors and infinite distance.

Thus we can find the optimal strategy for the initial peg solitaire configuration by forming $`G`$, computing its shortest path, using the location of the edges from $`(t,i)`$ to $`(s,i)`$ to partition the configuration into one-peg subconfigurations, and applying Corollary 1 to each subconfiguration. Since $`|G|=𝒪(n)`$, this algorithm runs in linear time.

In contrast to these results, Uehara and Iwata  showed that in two or more dimensions peg solitaire is NP-complete. However, the complexity of finding the minimum number of pegs to which a $`k\times n`$ configuration can be reduced, for bounded $`k>2`$, remains open.

Finally, we note that Ravikumar has proposed an impartial two-player game, in which players take turns making peg solitaire moves, and whoever is left without a move loses. It is tempting to think that this two-player game might be PSPACE-complete in two or more dimensions, and polynomial-time solvable in one.

### Acknowledgements.

We thank Aviezri Fraenkel, Michael Lachmann, Molly Rose, B. Sivakumar, and Spootie the Cat for helpful conversations, and the organizers of the 2000 MSRI Workshop on Combinatorial Games.


text:
# Statistical Analysis of Genealogical Trees for Polygamic Species.

## Abstract

Repetitions within a given genealogical tree provides some information about the degree of consanguineity of a population. They can be analyzed with techniques usually employed in statistical physics when dealing with fixed point transformations. In particular we show that the tree features strongly depend on the fractions of males and females in the population, and also on the offspring probability distribution. We check different possibilities, some of them relevant to human groups, and compare them with simulations.

One of the main problems encountered in the efforts to preserve species from extinction is genetic diversity. Indeed, besides environmental threats to the welfare of a species, a less obvious but nonetheless extremely important issue is related to the largeness of the genetic pool from which the genes of an individual are taken. Such a problem is related to the degree of consanguinity within the population: the more relatives mate among themselves, the more reduced is the genetic diversity of the population. There are examples in the wilderness of species with a relatively small genetic variety: from molecular biology it is known that cheetahs, for example, show a high degree of consanguinity, probably due to some bottleneck in the number of individuals in their population some ten thousands of years ago; in human societies, it is well known that high rank aristocrats in Europe kept marrying only among themselves. As a consequence, the appearance of a hemophiliac individual spread the genetic disease all over the reigning houses of Europe. This example sheds light on the relevance of the genetic diversity of a population for its conservation: species with a small genetic pool are weaker against genetic diseases. The above examples show that genetic redundancy can come as a consequence of a reduced population.

In this Letter we address the same problem from a different (but we believe complementary) standpoint: we are interested in the genealogical trees of individuals of species where the male-to-female ratio is not $`1`$ as in humans (here we define this ratio taking into account only males and females that are sexually mature). Among such examples we can name lions, sea lions and some antelopes, where each successfully reproducing male mates with more than one female (similar arguments could also be applied to polygamic human groups). Extreme cases are insects like bees and termites, where for every reproductive female (queen) there are very many males.

We measure the genetic redundancy in the gene pool of an individual by measuring the number of times that one of its ancestors many generations in the past appears more than once in its genealogical tree. Indeed, if no relatives would mate among themselves then, since every individual has a mother and a father, it would have $`2^g`$ ancestors $`g`$ generations in the past, half of them males and half of them females. Each of them would appear only once in the genealogical tree of their present descendents. Going back some tens of generations in the past, the number of ancestors would largely exceed the population itself. The only way out from this paradox is to assume that relatives indeed mate among themselves. As a consequence some individuals appear more than once in the genealogical tree of their descendents (that is, more than one branch of the tree had origin from such individuals), thus reducing the genetic pool from which their genes are taken.

We take a population of $`N`$ individuals, and we assume that it does not change in time. There is a fraction $`fN`$ of males and $`(1f)N`$ of females, ad also this fraction remains constant in time. Every male mates therefore, on the average, with $`1/f1`$ females. Here in general we make the (politically uncorrect) assumption that the fraction of males is less than $`1/2`$. Since in this model there is no difference between males and females, the opposite situation is obtained with a transformation $`f1f`$ (everything is symmetric with respect to $`f=1/2`$). We apply and extend the same scheme as developed in , generalizing it to the case of male fractions different from $`1/2`$.

Given an individual in the present generation, we are interested in the number of times its ancestors at a previous generation $`g`$ appear in the genealogical tree of that individual (at $`g=1`$ we find parents, at $`g=2`$ the grandparents, and so on). We therefore define $`m_r(g)`$ ($`f_r(g)`$) as the number of males (females) appearing $`r`$ times at generation $`g`$ in the genealogical tree of an individual at generation $`0`$, the present one.

The normalization of $`m_r(g)`$ and $`f_r(g)`$ implies that we can write

$$\underset{r=0}{\overset{\mathrm{}}{}}m_r(g)\mathrm{\Delta }r=fN,\underset{r=0}{\overset{\mathrm{}}{}}f_r(g)\mathrm{\Delta }r=(1f)N$$
(1)
where $`\mathrm{\Delta }r=1`$ trivially (but it is useful to write it explicitly for future rescalings). Since an individual at generation $`0`$ has $`2^{g1}`$ male ancestors (not necessarily distinct) at generations $`g`$ (and $`2^{g1}`$ female ancestors as well), we can also write

$$\underset{r=0}{\overset{\mathrm{}}{}}rm_r(g)\mathrm{\Delta }r=\underset{r=0}{\overset{\mathrm{}}{}}rf_r(g)\mathrm{\Delta }r=2^{g1}.$$
(2)
We define then the probabilities connected to $`m_r(g)`$ and $`f_r(g)`$. These are probabilities defined over the population at generation $`g`$. Therefore we have

$$M_r(g)=\frac{m_r(g)}{fN},F_r(g)=\frac{f_r(g)}{(1f)N}.$$
(3)
Using (3) we rewrite (1) as

$$\underset{r=0}{\overset{\mathrm{}}{}}M_r(g)\mathrm{\Delta }r=\underset{r=0}{\overset{\mathrm{}}{}}F_r(g)\mathrm{\Delta }r=1$$
(4)
and (2) as

$$\underset{r=0}{\overset{\mathrm{}}{}}rM_r(g)\mathrm{\Delta }r=\frac{2^{g1}}{fN},\underset{r=0}{\overset{\mathrm{}}{}}rF_r(g)\mathrm{\Delta }r=\frac{2^{g1}}{(1f)N}.$$
(5)
Finally we rescale $`r`$, $`F_r(g)`$ and $`M_r(g)`$ as

$`P_M(r,g)={\displaystyle \frac{2^{g1}}{fN}}M_r(g)`$ , $`P_F(r,g)={\displaystyle \frac{2^{g1}}{(1f)N}}F_r(g)`$ (6)
$`w_M(g)={\displaystyle \frac{fN}{2^{g1}}}r`$ , $`w_F(g)={\displaystyle \frac{(1f)N}{2^{g1}}}r.`$ (7)
With these definitions Eqs.(4) become

$$_0^{\mathrm{}}P_M(w_M,g)𝑑w_M=_0^{\mathrm{}}P_F(w_F,g)𝑑w_F=1$$
(8)
and Eqs.(5) become

$$_0^{\mathrm{}}w_MP_M(w_M,g)𝑑w_M=_0^{\mathrm{}}w_FP_F(w_F,g)𝑑w_F=1.$$
(9)
From (8) we see that $`P_M(w_M,g)`$ and $`P_F(w_F,g)`$ can be considered true probabilities. Next, we can write a system of equations for $`w_m(g)`$ and $`w_F(g)`$. A male $`i`$ at generation $`g+1`$ in the past has a number of repetitions that is given by the number of repetitions of his children at generation $`g`$. Therefore

$$r_{M,i}(g+1)=\underset{jsonofi}{}r_{M,j}(g)+\underset{jdaughterofi}{}r_{F,j}(g)$$
(10)
and analogously for females

$$r_{F,i}(g+1)=\underset{jsonofi}{}r_{M,j}(g)+\underset{jdaughterofi}{}r_{F,j}(g)$$
(11)
Dividing the first equation for $`2^{g1}/fN`$ we get

$$w_{M,i}(g+1)=\frac{1}{2}\underset{jsonofi}{}w_{M,j}(g)+\frac{f}{2(1f)}\underset{jdaughterofi}{}w_{F,j}(g).$$
(12)
Dividing (11) for $`2^{g1}/(1f)N`$ we get the analogous equation for females.

We assume a stable (on the average) population of $`N`$ individuals divided in two parts whose proportions are also (on the average) stable. Therefore the number of sons (daughters) that an individual can have has to obey well defined probability distributions. In our simulations we proceed backward in time, keeping the population fixed at $`N`$ and the male proportion fixed at $`f`$. Since we assign to every individual a couple of parents at random in the previous generation, the corresponding son/daughter probability distributions are binomials distributions. More precisely, the probability that a male has $`k`$ sons is

$$p_{mm}(k)=\left(\begin{array}{c}fN\\ k\end{array}\right)\left(\frac{1}{fN}\right)^k\left(1\frac{1}{fN}\right)^{fNk}$$
(13)
and that he has $`k`$ daughters is

$$p_{mf}(k)=\left(\begin{array}{c}(1f)N\\ k\end{array}\right)\left(\frac{1}{fN}\right)^k\left(1\frac{1}{fN}\right)^{(1f)Nk}$$
(14)
Analogous distributions can be written for $`p_{ff}(k)`$ and $`p_{fm}(k)`$.

We assume that the population is very large ($`N\mathrm{}`$) and that all the $`w`$’s are independent (this is verified in the limit of large $`N`$). In this limit the offpring probabilities become

$`p_{mm}(k,f)`$ $`=`$ $`p_{ff}(k,f)={\displaystyle \frac{e^1}{k!}}`$ (15)
$`p_{mf}(k,f)`$ $`=`$ $`p_{fm}(k,1f)={\displaystyle \frac{e^{(1f)/f}}{k!}}\left({\displaystyle \frac{1f}{f}}\right)^k.`$ (16)
In the case $`f=1/2`$ we recover the distributions used in .

Upon defining the generating functions

$`G_g(\lambda )`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}e^{\lambda w_M}P_M(w_M,g)𝑑w_M`$ (17)
$`H_g(\mu )`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}e^{\mu w_F}P_F(w_F,g)𝑑w_F`$ (18)
we find then that (12) become

$`G_{g+1}(\lambda )`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}p_{mm}(k)\left[G_g\left({\displaystyle \frac{\lambda }{2}}\right)\right]^kp_{mf}(j)\left[H_g\left({\displaystyle \frac{\lambda }{2}}{\displaystyle \frac{f}{1f}}\right)\right]^j`$ (19)
$`H_{g+1}(\mu )`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{\mathrm{}}{}}}{\displaystyle \underset{j=0}{\overset{\mathrm{}}{}}}p_{fm}(k)\left[G_g\left({\displaystyle \frac{\mu }{2}}{\displaystyle \frac{1f}{f}}\right)\right]^kp_{ff}(j)\left[H_g\left({\displaystyle \frac{\mu }{2}}\right)\right]^j`$ (20)
where also the equation for females has been written explicitely.

Substituing (16) in (20) we get, after some algebra,

$`G_{g+1}(\lambda )`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{f}}+G_g\left({\displaystyle \frac{\lambda }{2}}\right)+{\displaystyle \frac{1f}{f}}H_g\left({\displaystyle \frac{\lambda }{2}}{\displaystyle \frac{f}{1f}}\right)\right]`$ (21)
$`H_{g+1}(\mu )`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{1f}}+{\displaystyle \frac{f}{1f}}G_g\left({\displaystyle \frac{\mu }{2}}{\displaystyle \frac{1f}{f}}\right)+H_g\left({\displaystyle \frac{\mu }{2}}\right)\right]`$ (22)
These equations are clearly symmetric in $`f1f`$, since we do not make any distinction between males and females apart from the male proportion $`f`$.

Next, we analyze the stationary equations, $`g=\mathrm{}`$,

$`G(\lambda )`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{f}}+G\left({\displaystyle \frac{\lambda }{2}}\right)+{\displaystyle \frac{1f}{f}}H\left({\displaystyle \frac{\lambda }{2}}{\displaystyle \frac{f}{1f}}\right)\right]`$ (23)
$`H(\mu )`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{1f}}+{\displaystyle \frac{f}{1f}}G\left({\displaystyle \frac{\mu }{2}}{\displaystyle \frac{1f}{f}}\right)+H\left({\displaystyle \frac{\mu }{2}}\right)\right]`$ (24)
The probability that a male (a female) in the past does not appear in the genealogical tree of a given individual in the present generation is recovered sending $`\lambda ,\mu \mathrm{}`$ (by tauberian theorems, the limit $`\lambda ,\mu \mathrm{}`$ corresponds to the limit $`r_M,r_F=0`$). Therefore, upon calling $`G_0=G(\mathrm{})`$ and $`H_0=H(\mathrm{})`$ we have

$`G_0=\mathrm{exp}\left({\displaystyle \frac{1}{f}}+G_0+{\displaystyle \frac{1f}{f}}H_0\right)`$ (25)
$`H_0=\mathrm{exp}\left({\displaystyle \frac{1}{1f}}+{\displaystyle \frac{f}{1f}}G_0+H_0\right)`$ (26)
These equations can be solved numerically and the solution is shown in Fig.1(Left) (the results of the simulations agree with this solution up to the third significative digit).

Next, we expand (24) around the fixed point assuming that $`P_M(w_M)G_0\delta (w_M)+w_M^{\beta _M}`$ and $`P_F(w_F)H_0\delta (w_F)+w_F^{\beta _F}`$ for $`w_M,w_F0`$, which translates, by tauberian theorems, to

$$G(\lambda )=G_0+A_M\lambda ^{\beta _M1},H(\mu )=H_0+A_F\mu ^{\beta _F1}$$
(27)
for $`\lambda ,\mu \mathrm{}`$. Eqs.(24) then become

$`G_0\left[2^{\beta _M+1}+{\displaystyle \frac{A_F}{A_M}}\left(2{\displaystyle \frac{1f}{f}}\right)^{\beta _F+1}\lambda ^{\beta _M\beta _F}\right]`$ $`=`$ $`1`$ (28)
$`H_0\left[2^{\beta _F+1}+{\displaystyle \frac{A_M}{A_F}}\left(2{\displaystyle \frac{f}{1f}}\right)^{\beta _M+1}\mu ^{\beta _F\beta _M}\right]`$ $`=`$ $`1`$ (29)
Eqs. (29) are well defined only if $`\beta _M=\beta _F=\beta `$, and therefore we get, after some algebra,

$$2^{\beta +1}(H_0+G_0)=1$$
(30)
from which we can calculate the exponent $`\beta `$ as a function of $`f`$, shown in Fig.1(Right).

From (26) it is also possible to get the analytic behavior of $`H_0`$, $`G_0`$ and $`\beta `$ close to $`f=0`$:

$$G_0e^{\sqrt{\frac{2}{f}}},H_01\sqrt{2f},\beta 1+\frac{\sqrt{2}}{\mathrm{ln}2}f^{\frac{1}{2}}$$
(31)
As an exemple of distributions, in Fig.2 we show $`M_r(g)`$ and $`F_r(g)`$, (3), and in the inset their rescaled counterpart according to (7), for $`f=1/16`$. The exponent $`\beta `$ is negative, as from our analytical calculations. The delta function for $`r=0`$ has been omitted for scale reasons.

The dependence of $`\beta `$ from $`f`$ shows that such an exponent is highly nonuniversal and that it is extremely sensitive to the explicit form of the distributions (16). This becomes important when looking at real data. In the thirties Lotka fitted the probability of a man to have $`k`$ sons in the United States by a geometric distribution $`p_k=b_{mm}c_{mm}^{k1}`$ for $`k0`$, and $`p_0=d_{mm}`$, with $`c_{mm}=0.5893`$, $`d_{mm}=0.4825`$ and $`b_{mm}`$ chosen for normalization. Clearly, such a distribution is not a Poisson distribution as used above. Moreover it would give a rate of increase in the population of $`N_g/N_{g+1}=1.26`$.

Since in the definition of $`P`$ and $`w`$ in (7) depend on $`g`$, the particular value of $`N_g`$ can be explicitly incorporated in it. The left hand side of (12) is now multiplied by $`N_g/N_{g+1}`$. The probabilities for a male to be son of a male and a female to be daughter of a female will be those of Lotka, and the other ones can be evaluated by maintaining the fraction of males and females in the population constant, which translates in the constraints:

$$\frac{1d_{mf}}{1c_{mf}}=\frac{1f}{f}\frac{N_g}{N_{g+1}},\frac{1d_{fm}}{1c_{fm}}=\frac{f}{1f}\frac{N_g}{N_{g+1}}.$$
(32)
We can then rewrite (22) as

$`{\displaystyle \frac{N_g}{N_{g+1}}}G_{g+1}(\lambda )`$ $`=`$ $`\left(d_{mm}+{\displaystyle \frac{(1c_{mm})(1d_{mm})G_g(\lambda /2)}{1c_{mm}G_g(\lambda /2)}}\right)\left(d_{mf}+{\displaystyle \frac{(1c_{mf})(1d_{mf})H_g(\lambda /2)}{1c_{mf}H_g(\lambda /2)}}\right)`$ (33)
$`{\displaystyle \frac{N_g}{N_{g+1}}}H_{g+1}(\mu )`$ $`=`$ $`\left(d_{fm}+{\displaystyle \frac{(1c_{fm})(1d_{fm})G_g(\mu /2)}{1c_{fm}G_g(\mu /2)}}\right)\left(d_{ff}+{\displaystyle \frac{(1c_{ff})(1d_{ff})H_g(\mu /2)}{1c_{ff}H_g(\mu /2)}}\right)`$ (34)
Here we examine two different cases. First, we take $`f=1/2`$ and all $`c`$s and $`d`$s as from Lotka. We find that the probability $`G_0=H_0=0.231`$, different from the one obtained with Poisson distributions. Then we impose that the population size remains constant, $`N_g=N_{g+1}`$, but allow for different male fractions. Moreover, for simplicity, we choose $`d=1c`$ for the four probability distributions, in such a way that they become genuine geometric distributions: $`p_{mm}(k)=p_{ff}(k)=1/2^{k+1}`$, $`p_{mf}(k)=f(1f)^k`$ and $`p_{fm}(k)=(1f)f^k`$. The results for $`G_0`$ and $`H_0`$ are shown also in Fig.1(Left). The exponent $`\beta `$ is shown in Fig.1(Right). $`G_0`$ and $`H_0`$ approach their limit for $`f0`$ as $`f^{1/2}`$. In particular, the values for $`f=1/2`$ are clearly different from the ones with Poisson distributions. We find therefore that neither $`G_0`$ and $`H_0`$, nor $`\beta `$ are universal, although their behavior with respect to $`f`$ does not, qualitatively, depend on the details of the chosen offspring distribution. Actually, the relevance of the distribution to be used is hardly overestimated: one should take distributions obtained from the analysis of real data, in order to draw more detailed conclusions.

The present results show that, besides bottlenecks in the population size, there may be other factors affecting the largeness of the genetic pool from which the genes of an individual are taken. Indeed, for species with a very low value of $`f`$ we find that most females do not contribute to the genes of an individual in the present generation, whereas most males (who are anyway a little fraction $`f`$ of the entire population) do. As an extreme case (and exchanging males with females), in the absence of interbreeding between different hives, a single bee queen gives its genes to all subsequent generations. Some genetic mutation will become rapidly a genetic trait of the whole progeny. In case of bad mutations, they could well wipe out the whole family line. Although not dangerous per se, since bees and alike are extremely numerous, such a feature can make the species more sensitive to population size fluctuations.

In conclusion, we have generalized and analyzed the model proposed in  to the realistic case of species and human groups with male-to-female mating ratios different from $`1`$. Our results point out that the genes of an individual are taken from a pool whose largeness strongly depends on the male-to-female ratio, with important consequences when the population size strongly fluctuates. We are currently investigating the coupling effects between these different factors. Yet our results, although qualitatively of a general applicability, clearly show that quantitative estimates can only come when the analytical treatment is implemented with field data, since, as it is evident from Figs.1(Left) and 1(Right), different offspring probability distributions give rise to different quantitative results. This is a highly non-universal problem.

We thank F. Guinea for useful comments and discussions. P. De Los Rios thanks the Instituto de Ciencia de Materiales in Madrid, where this work was begun, for its kind hospitality. This work has been partially supported by the European Network contract FMRXCT980183.


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# Classical trajectories compatible with quantum mechanics

## I INTRODUCTION

Suppose that a non-relativistic particle is allowed to move in one dimension $`x`$ (for simplicity) under the influence of a potential function. There are two general perspectives for analyzing the motion of the particle: (i) that of classical physics, according to which the particle obeys Newtonian mechanics, or (ii) that of quantum mechanics, according to which the particle obeys, e.g., the Schroedinger wave equation (SWE). (We ignore for the moment the combined quantum-classical approach of D. Bohm .) Here we will show that there is, in fact, a third, intermediary viewpoint, whereby the particle obeys to a degree both classical mechanics and quantum mechanics.

Let the particle have mass $`m`$, and be moving under the influence of a stationary potential $`V(x)`$. Later on we treat the more general, time-dependent case. The particle is otherwise not perturbed by any exterior effect including measurement. Assuming non-relativistic speeds, the kinematics of the particle obey the Schroedinger wave equation (SWE). In the presence of our stationary potential this has as its solution a separable wave amplitude function $`\mathrm{\Psi }(x,t)=\psi (x)\mathrm{exp}(iEt/\mathrm{})`$, where $`\psi (x)`$ is any complex amplitude function, $`E`$ is a definite energy value, and $`\mathrm{}`$ stands for Planck’s constant (divided by $`2\pi `$).

The probability density function (PDF) for the separable case loses the time-dependence, $`p_X(x)=|\psi (x)|^2`$. This means that, from a quantum viewpoint, the coordinate $`x`$ is independent of the time $`t`$. The same PDF on position is obeyed at all times. However, from a classical viewpoint, in the presence of the potential $`V(x)`$ a particle has a definite trajectory $`x(t)`$ through time. To what extent can the quantum result and the classical viewpoint agree?

## II TRAJECTORY

Regard the particle as having a definite trajectory $`x(t)`$ during a time interval $`(0,T)`$. Consider the histogram of position values $`x`$ taken on by the particle during that time. The main point of this paper is that for any quantum solution $`\psi (x)`$ there is a classical trajectory $`x(t)`$ whose density of position values $`x`$ over the given time interval exactly coincides with $`|\psi (x)|^2`$. This trajectory is found next.

Central to the calculation is the following effect. Among the four coordinates of space-time, each time coordinate value is unique in occurring once and only once over a given interval $`(0,T)`$. That is, the time values as coordinates occur predictably in sequence one after the other. Of course, such a sequence of numbers has a density $`p_T(t)`$ of values that is uniform or flat over any time interval $`(0,T),`$

$$p_T(t)=1/T,t=(0,T).$$
(1)
This expression is unusual in expressing a density of deterministic events. However, whether deterministic or random, the density still must obey Jacobian transformation theory .

One defining property of a trajectory $`x(t)`$ is that to each time value $`t`$ in the sequence there is one and only one $`x`$ value. Consider the density $`p_X(x)`$ of all $`x`$ values traveled over by the particle during the given time interval. Note that like $`p_T(t),`$ density $`p_X(x)`$ is not a probability density, since the events $`x`$ follow a deterministic trajectory $`x(t)`$. Since to each $`t`$ value there is only one $`x`$ value the frequency of occurrence of each $`x`$ (actually, of each interval $`(x,x+dx)`$) equals that of its corresponding value $`t`$. Or, in terms of density functions ,

$$p_X(x)|dx|=p_T(t)|dt|,$$
(2)
where $`x=x(t)`$. Combining this with Eq. (1) gives

$$|\frac{dx}{dt}|=\frac{1}{Tp_X(x)}.$$
(3)
The absolute value sign for dx/dt defines two possible scenarios: either x increases with t during the entire time period T or x decreases with t during the period. (Note that for reasons of continuity in x, the choice of sign cannot be changed during the period.) By either sign, x is a monotonic function of t. Thus the particle makes one pass through the system in either the positive x direction or in the negative x direction. However, since the trajectory obeys the required PDF $`p_X(x)`$ via Eq. (3) for either direction (either choice of sign for dx/dt), the particle can alternatively obey periodic motion. This would be for any number of such periods T, i.e., any number of traversals to and fro through the system. At each such traversal the particle position would be a monotonic function of the time (as above) as defined by Eq. (3). In summary, the particle motion can be either single pass, or oscillatory.

To be definite, consider the monotonically increasing orbit dx(t)/dt $`0.`$ Let the density $`p_X(x)`$ of deterministic events $`x`$ equal the density $`|\psi (x)|^2`$ of random events $`x`$ as given by the SWE. Then Eq. (3) shows that the probabilistic density $`|\psi (x)|^2`$ is achieved by a classical trajectory $`x(t)`$ obeying

$$\frac{dx}{dt}=\frac{1}{T|\psi (x)|^2}.$$
(4)
This shows a possible problem at positions x for which $`|\psi (x)|^2`$ $`0.`$ The required particle speed dx/dt becomes so large as to exceed c, the speed of light. However, Eq. (4) also shows an inverse dependence upon the period T, and T has yet to be chosen. Obviously, if $`|\psi (x)|^2`$ is small but not zero the problem can be avoided by making T large enough. Furthermore, even at x for which $`|\psi (x)|^2`$ $`=0,`$ with T chosen large enough the unphysical region x within which relativity is violated can be made as narrow (as improbable) as is required.

Continuing with the analysis, Eq. (4) is straightforward to integrate, giving

$$𝑑x|\psi (x)|^2=(tt_0)/T,$$
(5)
with $`t_0=const`$. The integral is an indefinite one, resulting in a well-defined function $`g(x)`$ corresponding to the given $`\psi (x).`$ Hence Eq. (5) expresses the trajectory as

$$g(x)=(tt_0)/T,$$
(6)
which has to be inverted to obtain the desired trajectory

$$x=g^1[(tt_0)/T]x(t).$$
(7)
This is the formal solution to the problem.

There are two distinct interpretations to this result. As it was derived, it gives the trajectory of $`x`$ values that, if binned deterministically at each time $`t`$ over the interval $`(0,T)`$ will give the required PDF $`|\psi (x)|^2`$. But alternatively, the binned trajectory must also give the required PDF if it is sampled with uniform randomness (1) in time. The analysis is blind to whether $`t`$ is random or deterministic, since it depends only upon the Jacobian transformation of density functions, and Jacobian theory holds irrespective of whether the densities define deterministic or random variables.

The time $`t_0`$ is of interest. With the function $`g(x)`$ fixed for a given problem, Eq. (6) shows that the value of $`t_0`$ is defined by knowledge of initial conditions. For example, if it is known that at $`t=0`$ position $`x`$ has a certain value $`x_0`$ then $`t_0`$ is fixed as $`Tg(x_0)`$. Therefore, to know $`t_0`$ requires observation of the particle at a particular space-time point. But if this observation were made the particle’s state would be changed, and it would no longer have the amplitude $`\psi (x)`$ that has been presumed in the calculation. In fact the state would be a function of time, contrary to the stationary state solution assumed.

Therefore the time $`t_0`$ remains an unknown constant of the theory, and the trajectory (7) suffers an unknown time displacement. We conclude that the trajectory is deterministic only up to an unknown additive constant. Another way of saying this is that the trajectory is one of a family of curves $`x(t;t_0)`$ where each member of the family is defined by a particular value of $`t_0`$. This defines, in fact, a stochastic process . We spoke at the beginning of defining a form of mechanics that is intermediary between classical and quantum mechanics. It is this stochastic process that is intermediary between the two, preserving aspects of each. Thus, the trajectory $`x(t)`$ is deterministic, or classical, to an additive random constant, and it is the latter that gives a quantum aspect to the trajectory.

It may be noted that any member of the ensemble of trajectories satisfies the required PDF $`|\psi (x)|^2`$. Eq. (4) depends upon the value of $`dx/dt`$ which, according to Eq. (6), is independent of the unknowable $`t_0`$.

## III CASE OF A PARTICLE IN A BOX

To find a particular trajectory we have to specify a particular case $`\psi (x)`$. Consider the case where the particle is in a box of length $`L`$. With the origin for coordinate $`x`$ at the center of the box, the wave function in the ground state obeys

$$\psi (x)=\sqrt{2/L}\mathrm{cos}(\pi x/L),L/2xL/2.$$
(8)
Using this in Eq. (5) gives a solution

$$x/L+(2\pi )^1\mathrm{sin}(2\pi x/L)=(tt_0)/T.$$
(9)
This defines the trajectory $`x(t)`$ of the particle for positions $`x`$ over the interval $`(0,L)`$. Since the left-hand side of Eq. (9) is transcendental in $`x`$, the equation must be inverted numerically for the desired trajectory $`x(t)`$. A particle that traverses the given curve $`x(t)`$ back and forth any integral number of times will satisfy the required PDF $`(2/L)\mathrm{cos}^2(\pi x/L).`$ As a check on the theory, the $`x`$-values along the curve can be ‘binned’ to form a histogram $`p(x)`$ of occurrences $`x`$. One would then ascertain that it matches the theoretical answer $`(2/L)\mathrm{cos}^2(\pi x/L)`$ to the accuracy defined by the finite bin size We have done this taking $`\mathrm{\Delta }x=0.02`$. The results are not displayed for the sake of saving space.

## IV NON-STATIONARY PROBLEMS

We may now generalize the approach. Suppose that the potential $`V`$ is now more generally a function of both $`x`$ and $`t`$. Then the SWE solution is of a generally non-separable form $`\mathrm{\Psi }(x,t)`$. Accordingly, the particle has a PDF on position that varies with the time,

$$|\mathrm{\Psi }(x,t)|^2=p(x|t).$$
(10)
The vertical solidus on the right-hand side emphasizes that this PDF on random values $`x`$ is at (conditional upon) a fixed, nonrandom value $`t`$ of time. May the particle still be describable by a trajectory $`x(t)`$?

The random variable in the amplitude function $`\mathrm{\Psi }(x,t)`$ is $`x`$. The time $`t`$ is assumed known. However, in practice, a particle is tracked with uniform randomness (1) in time. The net PDF on $`x`$ over all such time values obeys the partition law of statistics ,

$$p_X(x)=𝑑tp(x|t)p_T(t).$$
(11)
Then by Eqs. (1) and (10), Eq. (11) becomes

$$p_X(x)=\frac{1}{T}𝑑t|\mathrm{\Psi }(x,t)|^2.$$
(12)
This is now used in place of quantity $`|\psi (x)|^2`$ in Eq. (5) to get the required trajectory. If the positions $`x`$ along this trajectory are sampled with uniform randomness in time, the resulting histogram will obey $`p_X(x)`$ given by Eq. (12).

## V MOMENTUM AND PHASE VALUES

We now turn to the question of the particle’s momentum values $`\mu `$ (note: this is not the usual notation $`p,`$ because of previous use of $`p`$ to denote probabilities.) The momentum of a purely classical particle obeys $`\mu =mv=mdx/dt`$, suggesting that a momentum ‘trajectory’ $`\mu (t)`$ could be found by simply differentiating $`d/dt`$ the previously found trajectory $`x(t)`$. However, in general this gives an ill-defined answer for the corresponding PDF on momentum (or velocity $`v`$), as shown next. Corresponding to Eq. (2) is

$$p_V(v)|dv|=p_T(t)|dt|,or,p_V(v)=\frac{1}{T|dv/dt|},$$
(13)
after use of Eq. (1). Next, by Eq. (4)

$$\frac{dv}{dt}=\frac{d}{dt}(\frac{dx}{dt})=\frac{d}{dt}(\frac{1}{T|\psi (x)|^2})=\frac{2}{T|\psi (x)|^3}\frac{d|\psi |}{dx}v,$$
(14)
after use of the chain rule of differentiation. Substitution in Eq. (13) gives

$$p_V(v)=\frac{|\psi |^3}{2|v(d|\psi |/dx)|}.$$
(15)
Although values of $`\psi (x)`$ are continuous, its derivatives $`d\psi /dx`$ can be non-continuous . At such velocities a PDF $`p_V(v)`$ computed via Eq. (15) would be ill-defined. Hence, this approach to finding $`p_V(v)`$ is untenable in general. Likewise the corresponding PDF on momentum by Jacobian transformation to $`\mu =mv`$ would be ill-defined in general. Hence, the momentum of our particle cannot be obtained as the mass times velocity. This leaves the momentum representation of Heisenberg ($`\mu i\mathrm{}grad`$) as the alternative. Thus, the probability amplitude $`\mathrm{\Phi }(\mu )`$ on momentum values $`\mu `$ is the usual Fourier transform of $`\mathrm{\Psi }(x,t)`$. This Fourier relation implies that the particle obeys the Heisenberg uncertainty principle, as it must.

In summary, this approach to particle dynamics defines the particle positions by the stochastic process $`x(t;t_0)`$, where trajectory $`x(t)`$ is deterministic up to the unknown constant $`t_0`$, and defines momentum values by working in the usual Heisenberg representation $`momentumi\mathrm{}grad`$.

A general wave amplitude $`\psi (x)`$ is complex and hence has a finite phase function. How is phase information reflected in the trajectory x(t)? The answer is that it is NOT, since by Eq. (4) $`dx/dt`$ is blind to the phase function. Instead, the phase function follows from the Heisenberg representation for momentum, by taking in the usual way the inverse Fourier-transform of the probability amplitude $`\mathrm{\Phi }(\mu )`$ on momentum.

## VI Effective potential $`\overline{V}`$

The trajectory x(t),  which we constructed to obey the SWE, can also be made to obey Newton’s laws, through the effect of an appropriate potential. Newton’s second law is

$$F=\frac{\overline{V}}{x}=m\frac{dv}{dt},$$
(16)
where $`\overline{V}`$ is to be an effective potential function. The acceleration $`dv/dt`$ borne of the SWE obeys Eq. (14). Using this in Eq. (16) forces the acceleration to be responding as well to a classical force F. Using as well Eq. (4) for v, gives

$$\frac{d\overline{V}}{dx}=\frac{2m}{T^2|\psi |^5}\frac{d|\psi |}{dx},\psi =\psi (x).$$
(17)
This equation defines the effective potential function $`\overline{V}`$ $`(x)`$that must be present for the particle to simultaneously obey the SWE and Newton’s second law. It may be integrated directly, to give

$$\overline{V}(x)=\frac{m}{2T^2}|\psi (x)|^4.$$
(18)
The dependence upon $`\psi `$ makes sense classically in that where the particle tends not to be, i.e. at x for which $`\psi (x)`$ is low, the potential tends to be high, acting as a classical barrier. Also, the negative sign means that the effective potential is always attractive, i.e., tending to produce periodic motion.

In the more general case (Sec. IV) of a time-dependent potential V(x,t), by similar steps the effective potential obeys

$$\overline{V}(x)=\frac{m}{2T^2}p_X(x)^2.$$
(19)
## VII Discussion

The semi-classical theory of D. Bohm  attempts, as here, to ascribe a classical trajectory to a particle that a priori obeys the SWE. As with our approach, the particle trajectory by the Bohm approach is only knowable to an unknown, additive constant. However, Bohm’s theory differs from ours both operationally (mathematically) and in its physical assumptions. Operationally, instead of Eq. (4), the Bohm trajectory for a particle obeys

$$\frac{dx}{dt}=\frac{j(x,t)}{|\psi (x,t)|^2},$$
(20)
where $`j(x,t)`$ is the quantum mechanical probability current. Comparing Eqs. (4) and (20) shows that, unless the current is unity, the two approaches give different trajectories for the particle. In general the current obeys

$$j(x,t)=\frac{i\mathrm{}}{2m}[\mathrm{\Psi }\frac{\mathrm{\Psi }^{}}{x}\mathrm{\Psi }^{}\frac{\mathrm{\Psi }}{x}].$$
(21)
This rarely has the value unity. For example, if $`\mathrm{\Psi }`$ is of the separable form $`\mathrm{\Psi }(x,t)=(1/L)^{1/2}\mathrm{exp}(ikx)\mathrm{exp}(iEt/\mathrm{})`$, Eq. (21) gives a current $`j=\frac{\mathrm{}k}{mL}`$. Or, where $`\mathrm{\Psi }`$ is of the separable form $`\mathrm{\Psi }(x,t)=\psi (x)\mathrm{exp}(iEt/\mathrm{})`$ with $`\psi (x)`$ real, Eq. (21) gives a current $`j=0`$. Hence the Bohm trajectories usually differ from ours. As an example, in the particle-in-box case above, $`j=0`$ so that by Eq. (20) the Bohm trajectory obeys dx/dt = 0, i.e., the Bohm particle does not move in the box. As with our effective potential Eq. (18), the Bohm particle moves in the field of an effective potential function that depends upon the wave amplitude function. However the Bohm effective potential, commonly called the ”quantum potential”, is of an entirely different form .

The physical underpinning of the Bohm approach consists of four assumptions. We list these in the following, along with corresponding assumptions of our particle model:

(1) A SWE solution $`\mathrm{\Psi }(x,t)`$ is not a probability amplitude but, rather, a “field” analogous to the electromagnetic field. Our approach preserves the standard probabilistic nature of $`\mathrm{\Psi }(x,t)`$ .

(2) The phase $`S(x,t)`$ of $`\mathrm{\Psi }(x,t)`$ is also a classical Hamilton-Jacobi function of mechanics. We do not assume this. Instead, the classical trajectory follows from the special nature of the time coordinate as in Eqs. (1) and (2).

(3) The Bohm semi-classical particle is acted upon by both the given potential function $`V(x,t)`$ and an added potential function called the “quantum potential” where the latter is a function of $`\mathrm{\Psi }(x,t)`$ as well. Our classical particle is acted upon in parallel by the two potential functions of the problem. These are (i) a conventional potential function V(x) that defines, via the SWE, the probability field $`|\psi (x)|^2`$ on position x that the particle is required to obey, and (ii) an effective potential function $`\overline{V}(x)`$ that constrains the classical particle to take the particular trajectory that satisfies the required $`|\psi (x)|^2`$. The effective potential function $`\overline{V}(x)`$ derives from the conventional potential function V(x) via solution $`|\psi (x)|^2`$ to the SWE and Eq. (18). In this manner the SWE and the conventional potential function V(x) fix both the statistics of the particle and its required trajectory (to an additive constant).

(4) The Bohm trajectory obeys Eq. (20), which is based upon the assumption that the quantum mechanical current j is also a classical current. As we found, this assumption in general implies trajectories that are different from ours. The assumption that the quantum mechanical current j is a classical current is difficult to verify, since the quantum mechanical current j is not a physical observable .

REFERENCES

1. D. Bohm, Phys. Rev. 85, 166 (1952); P. R. Holland, The quantum theory of motion (Cambridge University Press, Cambridge, England, 1993).

2. B.R.Frieden, Probability, Statistical Optics and Data Testing, 2nd. ed. (Springer-Verlag, New York, 1991) .

3. L.I. Schiff, Quantum Mechanics (McGraw-Hill, New York, 1955) .

4. D.Z. Albert, Quantum Mechanics and Experience (Harvard Univ. Press, 1992); also see “Bohm’s alternative to quantum mechanics”, Scientific American,  May, 1994, 58-67 for other references to Bohm’s work.