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# Angle-dependent normalization of neutron-proton differential cross sections

## I Introduction

A measurement of the differential cross section for elastic neutron-proton ($`np`$) scattering is notoriously difficult. It is even so difficult that almost none of the data sets measured at energies below the pion-production threshold is completely free of systematic flaws. In partial-wave analyses (PWA’s) of the $`np`$ scattering data  some of these flaws do not give rise to sizable systematic contributions to $`\chi ^2`$. Data sets with such minor flaws will not distort too much the statistics  in these PWA’s, and therefore such sets can be included in the database; examples are the LAMPF data  and the TRIUMF data . Some flaws, on the other hand, are so serious that their contribution to $`\chi ^2`$ dominates over the statistical contribution to the extent that the standard rules of statistics no longer apply. Consequently, such data sets  must be excluded from the databases used in PWA’s. This is of course an unfortunate and undesirable situation, especially in view of the waste of investment and effort involved in these experiments.

In this paper we present the “adnorm” method. This is a method of Angle-Dependent NORMalization  to treat certain systematically-flawed $`np`$ differential cross sections. This adnorm method is meant to be used only in good, energy-dependent PWA’s. We will show that the application of this method to certain data sets, which were previously unacceptable in energy-dependent PWA’s, can give impressive results. The values of $`\chi ^2`$ drop dramatically and can even become statistically acceptable. This implies that these data sets can, instead of discarding them, from now on be included in the $`np`$ database. The salvation of these systematically-flawed data sets  is a major accomplishment of our adnorm method.

In recent publications  we pointed out that the Uppsala data at 162 MeV  contain unexplained large systematic errors. Also some other $`np`$ differential cross section measurements appeared to have systematic errors similar (but not identical) to the Uppsala data; the Princeton  and Freiburg  data are prominent examples.

In the following we will compare the data in the standard way with the energy-dependent Nijmegen partial-wave analysis PWA93 . First of all we will establish that these flawed data sets have significant, smoothly angle-dependent, systematic errors. Since we have no explanation for these systematic experimental errors, we must simply accept the fact that certain data have such errors. We are surprised, however, that so many $`np`$ differential cross sections have such similar, angle-dependent systematic errors. Then we will apply the adnorm method to these data and demonstrate that this method can correct for some of such systematic errors. In order to save two sizable data sets, which became almost acceptable after using the adnorm method, we changed slightly the definition of an individual outlier. This has nothing to do with the adnorm method, but it is a measure taken in the same spirit: Try to be as frugal as possible with data sets, and do not omit them from the data base, unless absolutely necessary. In this way we obtain normalized data sets that are statistically acceptable and that can be included henceforth in the $`np`$ databases for PWA’s.

## II Normalization

An $`np`$ differential cross section $`\sigma (\theta ,expt)`$, consisting of $`N_{data}`$ data points, is called experimentally normalized, when for this data set the normalization has actually been measured. In that case, it has an experimental norm $`N(expt)=1.00`$ with a corresponding experimental error $`\delta N(expt)`$. This norm and error are included as a datum in the data base. For backward $`np`$ scattering this error is often of the order of 4% or larger. In the Nijmegen PWA’s we also determine for each data set the normalization $`N(pwa)`$ with an error $`\delta N(pwa)`$. This is a calculated normalization (for a discussion of these points see Ref. ). This error $`\delta N(pwa)`$ is in most cases less than 1% . When these two normalizations $`N(expt)`$ and $`N(pwa)`$ differ by more than 3 standard deviations (s.d.), we remove the experimental normalization and its error from the database. The data set is then “floated”, which means in practice that a very large normalization error is assigned to the data set. A data set is also floated when its normalization has not been measured at all. In the case of floated data the calculated normalization $`N(pwa)`$ is determined solely by the angular distribution. The number of degrees of freedom $`N_{df}`$ for a set with $`N_{data}`$ data points is then $`N_{df}=N_{data}1`$. For the determination of the calculated normalization $`N(pwa)`$ we have sacrificed one degree of freedom.

For seriously flawed data sets one cannot get a sufficiently low value of $`\chi ^2`$ by merely adjusting or floating the normalization. Such data sets are then omitted from the data base for PWA’s. The main point of this paper is the observation that for some of such unacceptable data sets we can introduce an angle-dependent normalization $`N(\theta )`$ in such a way that we essentially sacrifice two or more degrees of freedom to obtain significant drops in the value of $`\chi ^2`$. Such a sacrifice is unfortunately necessary in our attempt to save these data sets from being discarded otherwise.

In the adnorm method we have to make assumptions when we are trying to parametrize the angular dependence of the systematic errors. We first map the experimental angular interval $`[\theta _{min},\theta _{max}]`$ onto the interval $`[1,1]`$. This mapping can be done in many ways. We consider the two mappings

$$x=(\theta \theta _+)/\theta _{},$$
(1)
with $`\theta _\pm =(\theta _{max}\pm \theta _{min})/2`$, and

$$x=(\mathrm{cos}\theta z_+)/z_{},$$
(2)
with $`z_\pm =(\mathrm{cos}\theta _{max}\pm \mathrm{cos}\theta _{min})/2`$.

For the discrete set of $`N_{data}`$ data points $`x_i`$ $`(i=1,N_{data})`$ on the $`x`$-interval $`[1,1]`$ we define the inner product $`(x,y)=(1/N_{data})_ix_iy_i`$. This allows us to construct the polynomials $`S_n(x)=_{i=0}^na_ix^i`$, which are orthogonal with respect to this inner product and normalized such that $`a_n=1`$. Next we expand $`N(\theta )`$ in these orthogonal polynomials $`S_n(x)`$ on this discrete set of data points. We write

$$N(\theta )=N_0\left[1+f(\theta )\underset{n=1}{\overset{p}{}}c_nS_n(x)\right].$$
(3)
We allow for the introduction of an extra function $`f(\theta )`$. In practical cases we make the simplest choice $`f(\theta )1`$, but e.g. $`f(\theta )=1/\sigma (\theta ,pwa)`$ could also be a suitable choice. The expansion in orthogonal polynomials gives exactly the same $`N(\theta )`$ as a power series expansion up to the same power $`p`$. In the case of a power series expansion the coefficients $`N_0`$ and $`c_n`$ vary very much with the value of $`p`$, this is not the case anymore for an expansion in orthogonal polynomials. The normalization $`N_0`$ and the $`p`$ adnorm parameters $`c_n`$ ($`n=1,p`$) and their errors are determined by the least-squares method, where the data are compared with PWA93. In an actual PWA such data sets contribute with $`N_{df}=N_{data}(1+p)`$ degrees of freedom. The $`\chi ^2`$ that results when $`p`$ adnorm parameters are introduced is called $`\chi _p^2`$. When the standard angle-independent normalization (with zero adnorm parameters) is applied, the $`\chi ^2`$ is called $`\chi _0^2`$.

How many adnorm parameters $`c_n`$ do we have to introduce? The basic rule is that each introduced parameter should cause a significant drop in $`\chi ^2`$. We apply a 3 s.d. criterion: We introduce the parameter $`c_n`$ only when $`\chi _{n1}^2\chi _n^29`$. The parameter is then significant. This procedure of introducing additional parameters stops when no significant drop in $`\chi ^2`$ can be achieved anymore. When we end up with nonzero adnorm parameters, then we have shown that there are significant angle-dependent systematic errors present in the data, and we have explicitly parametrized these systematic errors. We have convinced ourselves that the final result of the parametrization is essentially independent of our specific assumptions. The renormalized differential cross sections, obtained by different ways of parametrization, were statistically practically the same.

## III Uppsala data

Let us see how this procedure works out for the Uppsala data  at $`T_L=162`$ MeV. At this neutron beam energy the $`np`$ differential cross section was measured in five overlapping angular regions. We ordered these sets by increasing neutron scattering angles and called the sets 1 to 5, where set 1 contains the data at the most forward angles and set 5 at the most backward angles . These data were then compared to PWA93. We removed the point at 93 from set 2 because it contributes more than 9 (3 s.d.) to $`\chi ^2`$.

In Table I we list the number of data $`N_{data}`$ in each set, the value $`\chi _0^2`$ obtained by just applying the standard angle-independent normalization (all adnorm parameters $`c_n0`$), the values of $`\chi _p^2(\theta )`$ obtained by applying the $`\theta `$-adnorm method of Eq. 1, and the values of $`\chi _p^2(z)`$ obtained by applying the $`z`$-adnorm method of Eq. 2.

From this Table I we see that the fits of the sets 1 and 4 improve significantly (a drop in $`\chi ^2`$ of much more than 9) when introducing only one adnorm parameter $`c_1`$. For sets 2 and 3 this is not true, and we will therefore take $`c_10`$ for these two sets. For set 5 we need two $`\theta `$-adnorm parameters, while only one $`z`$-adnorm parameter is necessary. After the $`\theta `$-adnorm method is applied with only one adnorm parameter the slope of the angle-dependent normalization $`N(\theta )`$ is called $`\alpha `$. The value of this slope can be found in the next-to-last row of Table I. This slope is used to compare systematic errors in different experiments. It is important to note that the slope for set 1 is positive, while the slopes for the sets 4 and 5 are negative. In the last row is given the variation in % of the normalization $`N(\theta )`$ in the case $`p`$=1 over the interval $`[\theta _{min},\theta _{max}]`$.

The combined data set has $`N_{data}=87`$ and when normalized in the standard angle-independent way (no adnorm parameters) we obtain $`\chi _0^2=243`$. This value is 12 s.d. higher than the expectation value $`\chi ^2=82(13)`$. The value for $`\chi ^2`$ drops to $`\chi ^2(\theta )=92`$ after introducing four $`\theta `$-adnorm parameters ($`c_1`$ for each of the sets 1 and 4, and $`c_1`$ and $`c_2`$ for set 5). With the $`z`$-adnorm method it drops to $`\chi ^2(z)=94`$ after introducing three $`z`$-adnorm parameters ($`c_1`$ for each of the sets 1, 4, and 5). The difference between the two adnorm methods is minor. In order to demonstrate this, we calculated the $`\chi ^2`$(dif) for the difference between the two differential cross sections obtained by the two adnorm methods. This $`\chi ^2`$(dif) is very low, only 0.8 for the total data set. The conclusion is that the large systematic errors of unknown origin present in the Uppsala data can be corrected for by using one of the adnorm methods. The drop of about 150 in $`\chi ^2`$ resulting from the introduction of only 3 $`z`$-adnorm (or 4 $`\theta `$-adnorm) parameters is impressive.

To present the data in a similar way as was done by the Uppsala group , we averaged the data in the overlap regions between the different sets. The difference $`\mathrm{\Delta }\sigma (\theta )=N(\theta )\sigma (\theta ,expt)\sigma (\theta ,pwa)`$ normalized in the various ways discussed, is presented in Fig. 1. In the top panel we show the data normalized in the Uppsala way and we get $`\chi ^2=393`$ for the 54 data points. In the middle panel of Fig. 1 we show the data normalized in the standard angle-independent way. This leads to $`\chi _0^2=135`$. In the bottom panel of Fig. 1 we present the data normalized with the $`\theta `$-adnorm method. We obtain then $`\chi ^2=59`$. From Fig. 1 one clearly sees the difference between the various ways the data have been normalized, and the enormous improvement obtained with the adnorm method.

## IV Freiburg data

Another place where the angle-dependent normalization procedure works impressively is the abundant Freiburg data . This data set consists of 4 different measurements (labeled expt I to expt IV) of $`np`$ differential cross sections, each at 20 beam energies between $`T_L=199.9`$ and 580.0 MeV, with a spacing of about 20 MeV. Because we compare with PWA93 we can only study data with energies less than 350 MeV, i.e. the 8 energies from 199.9 MeV to 340.0 MeV. Because of their too high individual contribution to $`\chi ^2`$ (more than 3 s.d.) we remove from the database the four data points (expt, $`T_L`$, $`\theta `$) = (II, 261.9 MeV, $`154.96^{}`$), (II, 300.2 MeV, $`148.34^{}`$), (II, 340.0 MeV, $`148.07^{}`$), and (III, 199.9 MeV, $`144.32^{}`$). For the total Freiburg data set we are left then with 859 data points.

In Table II we present the $`\chi _0^2/N_{data}`$ values for these 4 experiments at 8 energies after the standard angle-independent normalization. It is clear that $`\chi _0^2`$ for most of these 32 data sets is much too high. For the total data set of 859 points we find $`\chi _0^2=2139`$, which is 32 s.d. higher than the expectation value $`\chi _0^2=827(41)`$. Therefore, the total Freiburg data set would normally be discarded in PWA’s. However, we can try the adnorm method. The results of applying the $`\theta `$-adnorm method are presented in Table III. The first striking observation is the enormous drop in $`\chi ^2`$, from 2139 to 831, for the 859 data points. This drop was achieved by introducing next to the original 32 normalizations $`N_0`$ also 45 $`\theta `$-adnorm parameters. This implies on the average a drop of no less than 29 per adnorm parameter.

Looking at the four experiments separately, one sees that expt I and expt IV have values for $`\chi ^2(\theta )`$ that are smaller than their expectation value ($`0.8`$ s.d. and $`0.8`$ s.d.), that expt II has a $`\chi ^2(\theta )`$ that is 0.6 s.d. higher than its expectation value, and that expt III has a $`\chi ^2(\theta )`$ that is 3.3 s.d. higher than its expectation value. This, unfortunately, means that expt III would have to be excluded from the database for PWA93. For the remaining 647 points of expt I, II, and IV we expect $`\chi ^2=588(34)`$. We get $`\chi ^2(\theta )=571`$, which is an excellent result. We have checked explicitly that the $`\chi ^2`$-distribution of the renormalized Freiburg data is in very good agreement with the theoretical expectation . With the $`z`$-adnorm method we got similar results. Using 46 $`z`$-adnorm parameters and 32 normalizations we reached $`\chi ^2(z)`$= 832 for the 859 data points. Also in this case expt III should be omitted from the database for PWA’s. In Table III we also present the values of the slope $`\alpha `$ of the data at $`T_L=199.9`$ MeV after applying the $`\theta `$-adnorm method with $`p=1`$.

The conclusion is that the Freiburg data set, when compared to PWA93 using the adnorm method, has three statistically acceptable experiments and one, expt III, that is statistically not acceptable. However, one must realize that if this expt III were included in a new PWA, then it might possibly have a statistically acceptable value of $`\chi ^2`$.

We would like to point out that also this expt III can be saved, when we are willing to bend a little our rule for individual outliers. In the Introduction we already pointed out that this has nothing to do with the adnorm method, but only with our wish not to discard data unless it is absolutely necessary. When for expt III a 2.5 s.d. rule is used instead of a 3 s.d. rule, we must remove also the three data points (III, 199.9 MeV, $`133.95^{}`$), (III, 240.2 MeV, $`149.25^{}`$), and (III, 340.0 MeV, $`130.47^{}`$) as more than 2.5 s.d. outliers. In Table III the most right column of expt III contains the relevant information for this case. For the 209 data points left from expt III we have the expectation value $`\chi ^2=191(20)`$, and we find $`\chi ^2(\theta )=238`$, which is 2.3 s.d. higher than expected. Therefore, expt III is now also statistically acceptable. The conclusion is that the four Freiburg experiments, consisting of 856 datapoints, can be made statistically acceptable with the adnorm method and using the 2.5 s.d. rule for outliers.

## V Princeton data

Finally, we consider the Princeton data . This relatively old data set is generally not included in PWA’s and these data were e.g. also discarded in the final version of PWA93. We can, however, revisit these data with the adnorm method. The results are given in Table IV. After we have removed two data points ($`T_L=313`$ MeV, $`\theta =168.1^{}`$ and $`170.3^{}`$) as more than 3 s.d. outliers, the total set contains 156 data points, divided over 9 energies below $`T_L=350`$ MeV. When we normalize these data in the standard manner we get $`\chi _0^2=582`$. This is about 25 s.d. higher than the expectation value. Next we applied the $`\theta `$-adnorm method. This required 14 additional adnorm parameters. The expectation value is then $`\chi ^2=133(16)`$. We obtain $`\chi ^2(\theta )=195`$, which is still 3.9 s.d. too high. Therefore the Princeton data, unfortunately, cannot be saved by the adnorm method alone, despite the enormous improvement in $`\chi ^2`$ from 582 to 195, which amounts on the average to a drop of 28 per adnorm parameter. Using the $`z`$-adnorm method gives similar results.

However, when again we are willing to bend our rule for individual outliers a little, we can also save these data. According to the 2.5 s.d. rule the three data points ($`T_L`$, $`\theta `$) = (224 MeV, $`131.6^{}`$), (239 MeV, $`139.7^{}`$), and (257 MeV, $`178.6^{}`$) must also be omitted. There are then 153 data points left, which leads to the expectation value $`\chi ^2=130(16)`$. The second entries in Table IV give the relevant information for this case. We obtain $`\chi ^2(\theta )=169`$, which is 2.4 s.d. higher than expected.

## VI Discussion and Conclusions

About the adnorm method that we proposed here, the question could be raised: “Is it successful because it corrects for experimental errors, or perhaps because it corrects for unknown biases in the PWA’s?” We claim that we correct for unknown systematic experimental errors. To demonstrate this we defined the slope parameter $`\alpha =(1/N_0)(dN(\theta )/d\theta )`$ for the case $`p=1`$. Looking at the different problematic experiments in about the same angular region and at about the same energy, we note significant differences. In the backward direction $`[150^{},180^{}]`$ there are several experiments at about the same energy. These are the Uppsala set 5 at 162 MeV, the Freiburg expt’s I and II at 199.9 MeV, and the Princeton data at 182 MeV. The values of $`10^3\alpha `$ are $`4.9(5)`$, $`2.8(4)`$, $`2.5(4)`$, and $`0.7(4)`$, respectively. These slopes do not agree! The disagreement between the Uppsala sets 1 and 2 and the Freiburg expt IV is worse. Uppsala set 1 covers the angular region $`[73^{},107^{}]`$ and has $`\alpha =3.7(8)\mathrm{\hspace{0.25em}10}^3`$. Uppsala set 2 covers $`[89^{},129^{}]`$ and has $`\alpha =0.3(7)\mathrm{\hspace{0.25em}10}^3`$. The Freiburg expt IV covers $`[81^{},124^{}]`$ and has $`\alpha =1.9(3)\mathrm{\hspace{0.25em}10}^3`$, which is even of the opposite sign as the values for the Uppsala sets 1 and 2. When one wants to blame PWA93 for the discrepancy and claim that the Uppsala and Freiburg data are in agreement, then one should at least find the same values for $`\alpha `$. Because the $`\alpha `$ values for these experiments are significantly different, we can conclude that the Uppsala and Freiburg data are not in agreement with each other. It is then also clear that the discrepancies must be of experimental origin.

In conclusion, we have shown that many of the measurements of the $`np`$ differential cross section suffer from similar systematic errors. These errors are mostly so large that the corresponding data sets cannot be included in PWA’s. However, it turned out that these systematic errors have a smooth angular dependence, which can easily be parametrized. This allowed us then to use the adnorm method to correct for these systematic errors. In many cases this gave rise to impressive drops in the values of $`\chi ^2`$ for several of the seriously flawed data sets. Many of the data sets became statistically acceptable after application of the adnorm method and therefore can now be included in the data base for PWA’s. However, some of the sets required also a slight change in the definition of outlier to make them acceptable. In this manner, several $`np`$ data sets  can be saved from oblivion.

###### Acknowledgements.

We thank M.C.M. Rentmeester, Th.A. Rijken, and R.A. Bryan for helpful discussions. The research of R.G.E.T. was made possible by a fellowship of the Royal Netherlands Academy of Arts and Sciences (KNAW).


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# More on“Atomic motions in the crystalline Al50Cu35Ni15 alloy”

## Abstract

We refute recent claims that ultrafast atomic jumps as observed in quasicrystals (QC) could be called phasons in many crystalline alloys by pointing out that there is a genuine conceptual difference between the hopping dynamics in an imperfect crystal containing a substantial number of vacancies, and the hopping dynamics due to phason motion in QC.

63.50, 66.30F, 64.70R

Recently we reacted to comments by Dahlborg et al. on work of phason hopping in quasicrystals (QC). We want to remove a number of misunderstandings that could result from the reply of these authors and Dolinšek.

It would appear as though we made a severe mistake by calling the sample a B2-phase rather than a $`\tau `$-phase. In one of the references we cited  the authors stated: “The structure of this alloy is of B2, CsCl type with lattice parameter $`a=`$ 32.9 Å.” Admittedly, the term B2-phase is outdated and should be replaced by CsCl-type phase. It has usually a broad stability range and may contain also structural vacancies. If the vacancies order leading to superstructures then these phases are called $`\tau `$-phases. If the B2-phase is fully disordered it is usually called $`\beta `$-phase. But it should be clear that an emphasis on such formal aspects of the presentation of our argument cannot address its contents. The issue is not how we name the sample, nor what its precise structure is, but that it is biased in having a significant structural-vacancy concentration.

Besides a B2 and a $`\tau `$ phase the authors call their sample also a one-dimensional QC. If true, this would completely undermine the pertinence of their sample as a counter example of an alloy that displays ultrafast hopping but is not a QC. Let us note in this respect also a tacit change of issues between the original paper (ultrafast atomic jumps are a trivial feature of all Al-based ternary alloys) and the reply (atomic jumps can be called phasons in many other solids than QC). In other words, the authors now claim that the concept of a phason (and of a QC) could be generalized to such an extent that it becomes a trivial object. They base these claims on a geometrical construction by Steurer.

Steurer’s paper aims at introducing a 1-1 mapping between a QC and a crystal. Before analysing the argument we want to point out that this starting point is not completely rigorous since mathematically spoken it is impossible to have such a 1-1-mapping, as has been conclusively shown by Duneau and Oguey. To express that we concede to neglect this problem we could state

(admittedly quite vaguely) that we allow for “imperfect” crystalline structures, although in principle this should require detailed further specification. Steurer’s construction consists then in defining two mappings: one, which is 1-1, between a QC and an average structure (AS) which is a (highly imperfect) periodic lattice with a well-determined structural-vacancy concentration, and one between the CsCl-type structure and the AS structure. In the case of an icosahedral QC (the Amman tiling), the second mapping can be the identity mapping. While it is not wrong that this way a QC can be put into registry with a highly imperfect CsCl-type structure with a uniquely determined structural-vacancy concentration (e.g. $`(3\tau 4)/(3\tau +1)`$ for the construction starting from a decagonal QC), it must be clear that on the contrary the reversal of this is just not true: An (imperfect) periodic structure which is the AS of a incommensurately modulated structure (IMS) is not itself an IMS, as the authors suggest. (Every perfect periodic structure can be used as the basic or average lattice for the definition of some IMS, which of course does not imply that all periodic structures would be IMS). Moreover, none of the cubic structures that have a different structural-vacancy concentration than the magic one, can be “considered” as a QC according to the philosophy of the authors. If the vacancy concentration in the AlCuNi sample were to be the magic one, then it would even be more biased.

Let us now discuss the generalized “phasons” of the authors. The second mapping does not use hyperspace, such that the nature of the phasons invoked by the authors must be obtained by transposing the picture of the phasons under the first mapping, using the second mapping as a dictionary. But the first mapping represents the QC as an IMS. A shift of the superspace cut in an IMS produces small atomic displacements of the atoms around the average positions, but leaves the positions of the AS unchanged (see e.g. Fig. 1b of reference \[5a\]). The images of these average positions, which are the atoms of the cubic phase, will thus not move eighter. Hence vacancy diffusion on the CsCl-type structure does not correspond to phasons (as the phasons remain restricted to the QC and do not lead to any atomic motion on the cubic lattice at all), but to true vacancy diffusion on the AS (and the QC). Intersticialcies on the CsCl-type structure are not addressed in Steurer’s construction.

As the reader can inspect on the Figure just mentioned, for large enough amplitudes of the translation of the cut in hyperspace the phason modes of the QC in the IMS setting correspond to a gradual (continuous) shift of the position up to a maximal distance followed by a (discontinuous) jump back to the starting position. That is a meaningless concept in the lattice dynamics of a cubic structure (and a QC).

There are no or very few data on hopping rates in ternary alloys. The claims of the authors that ultrafast hopping would be common in ternary Al-based alloys can therefore not be refuted by showing ample experimental evidence from the literature that would prove it wrong. However, there is more to scientific credibility of a statement than the absence of experimental information contradicting it. The statement of the authors is based on the extrapolation of a single data point collected on a heavily biased sample. Pointing this out should be sufficient to dismiss the entire argument. There is no further charge of proof on our shoulders (Actori incumbit probatio). Still, for further appraisal of the unusual, ad hoc character of this extrapolation, we showed that the jump rates observed in QC are almost a factor of 100 faster than in all other metallic alloys studied, by citing the examples with the fastest hopping rates ever observed in any non-QC alloy. (That these examples just happened to be B2 phases is here not part of the argument). By replying to this that one cannot compare a ternary phase to a binary phase, the authors mispresent the issues.

They also mispresent the issues when they claim that we would be speculating about the temperature dependence. If they want to prove that the fast jumps observed in QC are a common feature of all Al-based ternary alloys, then their temperature dependence has to be of the assisted type like in QC. Else the situations in the two samples just do not compare. Exactly this crucial information about the temperature dependence in their own sample is missing, a deficiency we flagged. Elsewhere, the authors tacitly assume a normal temperature dependence in their sample, by invoking a thermal-vacancy mechanism for the fast jumps.

It is, perhaps, not much of a “surprise” that we take great discomfort with a 21-pages Comment on our work having been presented as a regular article, and the inordinate delay of what just should have been our right of reply and the final word in the discussion.

Finally, we would like to signal an erratum that sneeked into our paper due to editorial changes during the production process. We wanted to quote Hf and $`\beta `$-Ti “show hopping rates two orders of magnitude slower” rather than “show self-diffusion two orders of magnitude slower”.

We would like to thank Prof. W. Steurer for providing us with some accurate definitions, and a copy of a paper \[5a\] we were unable to get access to.


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# Untitled Document

Comment on Formulating and Generalizing Dirac’s, Proca’s, and Maxwell’s Equations with Biquaternions or Clifford Numbers111Dedicated to Freeman Dyson.

André Gsponer and Jean-Pierre Hurni222Independent Scientific Research Institute, Box 30, CH-1211 Geneva-12, Switzerland
e-mail: isri@vtx.ch

Published in Foundations of Physics Letters, Vol. 14, No. 1 (2001) 77–85.

*Many difficulties of interpretation met by contemporary researchers attempting to recast or generalize Dirac’s, Proca’s, or Maxwell’s theories using biquaternions or Clifford numbers have been encountered long ago by a number of physicists including Lanczos, Proca, and Einstein. In the modern approach initiated by Gürsey, these difficulties are solved by recognizing that most generalizations lead to theories describing superpositions of particles of different intrinsic spin and isospin, so that the correct interpretation emerges from the requirement of full Poincaré covariance, including space and time reversal, as well as of reversion and gauge invariance. For instance, the doubling of the number of solutions implied by the simplest generalization of Dirac’s equation (i.e., Lanczos’s equation) can be interpreted as *isospin*. In this approach, biquaternions and Clifford numbers become powerful opportunities to formulate the Standard Model of elementary particles, as well as many of its possible generalizations, in very elegant and compact ways.*

1. INTRODUCTION

In the years that followed its discovery, various attempts were made to recast Dirac’s equation (1928) using biquaternions (Lanczos, 1929), spinors (van der Waerden, 1929), Clifford numbers (Proca, 1930), semivectors (Einstein, 1932), etc.(1,2) However, only the spinor and the original matrix formulations gained a wide acceptance.

In the 1950-1960’s, biquaternion and Clifford number formulations were revived by Gürsey, Riesz, and Hestenes. At the same time, the Dirac matrix formulation was more and more often taught with reference to the $`C\mathrm{}_{4,1}`$ Clifford algebra, and Feynman made some effort to introduce the two-component formalism based on the Pauli algebra $`M_2()C\mathrm{}_{3,0}`$ .

Yet it is only in last two decades of the XXth century that a growing number of researchers became truly interested in the formulation of fundamental physics using Clifford numbers. While such formulations have several advantages (e.g., they generally lead to neat expressions that are often very suggestive), they are affected by a number of difficulties, and they often lead to speculative interpretations, which tend to hinder their acceptance by the research community at large.

In particular, compared to ordinary matrix or tensor formulations, there are: *(i)* losses of information because Clifford numbers do not carry spinor or tensor indices to indicate their behavior under coordinate transformations, (ii) ambiguities due to the noncommutative nature of Clifford numbers (reversion333*Reversion* is also called *ordinal conjugation* in the quaternion literature because it consists of reversing the order of the factors in a monomial.) and to the residual arbitrariness in their definition,(2,3) *(iii)* difficulties of interpretation coming from the increase in the number of components when vectors or spinors are replaced by Clifford numbers,(4,5) *(iv)* similar difficulties when ordinary scalars such as the mass are promoted to Clifford numbers, and *(v)* more fundamental difficulties with the quantum interpretation when real Clifford algebras are used instead of complexified ones. The last difficulty comes from the well known fact that if *“i”* is not explicit in a quantum theory, it must be replaced by the definition of special rules(6) or by an equivalent antiunitary operator.(7)

In a short comment it is not possible to address all these difficulties in detail and to refer to the numerous papers in which they arise. (For a selection, see Ref.8–12 and references therein.) We therefore take a simple but non trivial case: the problem of formulating physical theories with biquaternions. Since these are just complexified quaternions, $`𝔹=C\mathrm{}_{3,0}`$ , there are no difficulties with the quantum interpretation. Nevertheless, biquaternions are suitable to show how the other difficulties can be circumvented, and how much the isomorphisms $`SU(2)/`$ and $`SL(2,)𝔹^{}/`$ are useful in making explicit the basic symmetries of the Standard Model.

2. DIRAC’S AND LANCZOS’S FUNDAMENTAL EQUATIONS

In the standard spinor or two-component formalisms, Dirac’s equation can be symbolically written as

$$\overline{}L=mRR=mL$$
$`(1)`$
where $`L`$ and $`R`$ are the left- and right-handed components of Dirac’s four-component bispinor, and $``$ the spinor or Pauli-matrix 4-gradient.

The simplest generalization of this equation (which amounts to replace the two-component spinors $`L`$ and $`R`$ by two biquaternions $`A`$ and $`B`$) has first been made by Lanczos:

$$\overline{}A=mBB=mA.$$
$`(2)`$
Here $`=[_{it},_\stackrel{}{x}]`$ is the quaternion 4-gradient operator, assumed to transform as a Lorentz 4-vector, i.e., $`^{}=^{}`$ with $`SL(2,)`$. This is *Lanczos’s fundamental equation*(4) which has the remarkable property of describing isospin doublets of spin $`\frac{1}{2}`$ particles, as well as spin 0 and spin 1 particles of either positive or negative space-reversal parity. The existence of isospin doublets is the consequence of the invariance of (2) under the $`SU(2)`$ gauge transformation $`AAG`$, $`BBG`$.

Dirac’s equation is obtained from (2) by making the superpositions:(4,13)

$$D_+=A\sigma +B^{}\overline{\sigma }D_{}=(A\overline{\sigma }B^{}\sigma )i\stackrel{}{e}_1$$
$`(3)`$
where $`\sigma `$ is a constant idempotent.444For definitiveness, we take $`\sigma =\frac{1}{2}(1+i\stackrel{}{e}_3)`$ where $`\stackrel{}{e}_3`$ is the third quaternion unit, and $`\stackrel{}{e}_1`$ for a typical unit vector perpendicular to $`\stackrel{}{e}_3`$. Both of these linearly independent superpositions satisfy the Klein-Gordon equation, $`\overline{}D=m^2D`$, as well as the Dirac-Lanczos equation:

$$\overline{}D=mD^{}i\stackrel{}{e}_3.$$
$`(4)`$
This equation is strictly equivalent to Dirac’s. It was first discovered by Lanczos$`^{(4,p.462)}`$ and later rediscovered by a number of people, in particular by Gürsey(14) and Hestenes.(15) In fact, in a form or another, it is the basic equation in any Clifford type formulation of Dirac’s theory.

The spin $`\frac{1}{2}`$ character of the superpositions (3) is enforced by the postmultiplication by $`\sigma `$, and in (4) by the factor $`ie_3`$, which imply that under a Lorentz transformation we have $`D^{}=D`$.

The correct interpretation of $`D_+`$ and $`D_{}`$ as an isospin doublet has been given for the first time by Gürsey.(13) This interpretation solves the “doubling” problem that puzzled Lanczos(4,2) and others more recently.(5,9)

3. PROCA’S EQUATION AND REVERSION

In order to obtain the spin 0 and spin 1 solutions of (2) one has to address a problem that is implicitly solved by making the superpositions (3) in the spin $`\frac{1}{2}`$ case. In effect, a fundamental concept used in expressing a physical theory in mathematical form is that any arbitrariness in the writing of the equations should have no consequences on the physics. While this principle is automatically taken care of in tensor calculus, in Clifford (or biquaternion) formulations it implies that the *order* of the factors (from left to right, or right to left) should have no consequencies.(2) In our case, the *reversed* version of equation (2)

$$A^{}\overline{}=mB^{}B^{}=mA^{}$$
$`(5)`$
could equally well have been taken as *Lanczos’s fundamental equation*. As shown in Ref.2, the concepts of reversion and space reversal are intimately connected: A requirement of definite reversion symmetry such as $`A^{}=\pm A`$ corresponds to selecting fields of even or odd parity. For example, if $`A=A^{}`$ is an odd parity 4-vector, it may represent Maxwell’s or Proca’s 4-potential. In this case, the corresponding reversion invariant field equation is obtained by combining (2) and (5):

$$\overline{}A=mBB+B^{}=2mA.$$
$`(6)`$
However, while the fields in (6) have a well defined parity, they are still mixtures of spin 0 and spin 1 states. This is because under a Lorentz transformation we have $`B^{}=^{}B^{}`$, so that the scalar part of $`B`$ is a Lorentz invariant, while the vector part transforms as a Lorentz 6-vector. Hence, the field $`[b,\stackrel{}{b}]=mB`$ can be covariantly decomposed into its scalar and vector parts. Introducing the corresponding decomposition of the potential $`A=A_0+A_1`$, the system (6) can be split into a spin 0 and a spin 1 part.

For the spin zero part, because $`A=A^{}=A_0`$, we see from (2,5) that $`b=b^{}`$. The system (6) reduces therefore to the following degenerate form of Lanczos’s equation:

$$\overline{}A_0=bb=m^2A_0.$$
$`(7)`$
For the spin one part, the resulting equation is just (6) with the field $`B`$ restricted to its vector part $`\stackrel{}{b}`$. This is Proca’s equation for a massive spin 1 particle:(2)

$$\overline{}A_1=\stackrel{}{b}\stackrel{}{b}+\stackrel{}{b}^{}=2m^2A_1.$$
$`(8)`$
The key assumption in going from Lanczos’s fundamental equation (2,5) to (7,8), i.e., the field equations for a massive “scalar” or “vector” particle, was that $`A`$ had to be a reversion symmetric 4-vector. If we had assumed the second fully covariant possibility, i.e., that $`A`$ was reversion anti-symmetric, we would have obtained the field equations for a pseudoscalar and respectively a pseudovector particle. The only difference is that in the second equation of (8) the plus sign has to be replaced by a minus sign. The fact that all equations for particles of spin 0 and 1 are degenerate cases of Lanczos’s equation has first been shown by Gürsey in his PhD thesis.$`^{(14,p.162)}`$

From Proca’s equation (8) it is now straightforward to derive the conserved current and the energy-momentum tensor(2)

$$C=A_{1}^{}{}_{}{}^{}\stackrel{}{b}+\stackrel{}{b}^{}A_1+(\mathrm{})^{}$$
$`(9)`$
$$8\pi T()=\stackrel{}{b}^{}()\stackrel{}{b}+m^2A_1^{}()A_1+(\mathrm{})^{}$$
$`(10)`$
where $`(\mathrm{})^{}`$ means that the expression has to be completed by adding the reversed version of the part on the left. If now the electromagnetic field is introduced, either by local gauge invariance or the minimum coupling substitution ansatz, the divergence of $`C`$ will still be zero, while the divergence of $`T()`$ will give the Lorentz force acting on the particle.

4. MAXWELL’S EQUATIONS AND REVERSION

Now that we have derived Proca’s equation, we can set $`m=0`$ in (8) to get the quaternionic field equation of a massless particle of spin 1. Including a source term with a current $`J`$ we obtain:

$$\overline{}\varphi =\stackrel{}{f}\stackrel{}{f}+\stackrel{}{f}^{}=8\pi J.$$
$`(11)`$
This is the correct form of Maxwell’s second set of equations in quaternion form. It explicitly implies that the source current $`J`$, and thus the potential $`\varphi `$, has to be reversion symmetric and shows that the electromagnetic field is the antisymmetric tensor:

$$2F()=()\stackrel{}{f}+\stackrel{}{f}^{}().$$
$`(12)`$
Contrary to the 6-vector $`\stackrel{}{f}=\stackrel{}{E}+i\stackrel{}{B}`$, this tensor is reversion symmetric, and therefore defines a physical field of odd parity. Its correct quaternionic form was first given by Kilmister.(3)

The usual quaternion form of Maxwell’s equations contains the electromagnetic 6-vector $`\stackrel{}{f}`$ instead of the electromagnetic tensor (12):

$$\overline{}\varphi =\stackrel{}{f}\stackrel{}{f}=4\pi J.$$
$`(13)`$
This equation is formally incorrect because it is not covariant under spatial reversal. However, provided that $`\varphi `$ and $`J`$ are postulated to be reversion invariant, it is in practice equivalent to (11) because from Lanczos’s fundamental system (2,5) we have then:

$$\stackrel{}{f}\stackrel{}{f}^{}=0.$$
$`(14)`$
This equation corresponds to the requirement that there should be no magnetic pole or current, i.e., to Maxwell’s first set of equations. Equation (14), however, is stronger than these equations in tensor form: it totally excludes any kind of magnetic source. Therefore, Lanczos’s equation and reversion invariance enforce a fundamental asymmetry that is strongly supported by experiment, i.e., the existence of point-like electric charges, and the absence of magnetic monopoles.(2)

5. EINSTEIN-MAYER’S EQUATION AND STANDARD MODEL

Looking at (2), an obvious generalization is to assume that the mass $`m`$ is a quaternionic parameter or field. This assumption leads to the concept of *hypercomplex mass*,(5) or to the *versatile Dirac equation*,(8) and was first made by Einstein and Mayer when they generalized Lanczos’s equation as follows:(16)

$$\overline{}A=BE^{}B=AE.$$
$`(15)`$
Depending on the suppositions made for $`E𝔹`$ one has a rich spectrum of solutions. As with Lanczos’s fundamental system, these have to be classified according to spin, parity, and mass. In field theory terminology, when $`E`$ is a spin 0 field such as the Higgs doublet of the Standard Model, (15) corresponds to a “Yukawa coupling” between $`E`$ and the $`A`$ and $`B`$ fields.(17)

A major difference with Lanczos’s fundamental system is that the second order equations for $`A`$ and $`B`$ are now eigenvalue equations for the mass. In other words, while all solutions of (2) have the same mass, (15) has pairs of solutions with different masses that Einstein and Mayer tried to interpret as electron-proton doublets.(16) In the modern perspective, two particularly interesting special cases emerge: First, when $`EE^{}=m^2`$ the spin $`\frac{1}{2}`$ isospin-doublet can be interpreted as a proton-neutron pair (because it is degenerate in mass). Then, if $`E`$ is assumed to be a pseudoscalar field, one is lead to the charge independent theory of low-energy nuclear forces.(18) Second, when $`E\overline{E}=0`$ the doublet can be interpreted as an electron-neutrino pair (because one of the two masses is necessarily zero). Then, assuming a parity violating $`SU(2)`$ local gauge such that $`AAG`$, $`BB`$ and $`EG^{}E`$, one gets a low-energy weak interaction phenomenology that is equivalent to that of the Standard Model.(19)

Other physical interpretations are possible when $`E`$ is derivatively coupled to the $`A`$ and $`B`$ fields. For example, if

$$E=A\overline{B}$$
$`(16)`$
the system (15-16) leads to double-periodic Petiau waves. These solutions interpolate between pure de Broglie waves and pure solitonic waves: a beautiful realization of the wave/particle duality of quantum mechanics.(20)

However, possibly the most interesting interpretation is to relate $`E`$ to the Higgs field of the Standard Model:

$$\overline{}E=\mu E+\lambda EE^{}E.$$
$`(17)`$
Since this equation is invariant in the quaternion gauge transformation $`EG^{}EH`$ it has the manifest $`SU(2)_L\times SU(2)_R`$ symmetry of the $`\sigma `$-Model which, after breaking the $`SU(2)_R`$ symmetry, leads to the renormalizable electroweak theory of the Standard Model.(21)

6. CONCLUSION

In this comment we have shown that the use of Clifford algebras, and more particularly of biquaternions, can lead to a satisfactory formulation of elementary particle physics which, in the Standard Model, is entirely based on fundamental fields of spin 0, $`\frac{1}{2}`$ and 1. In order to build such a formulation on solid ground, it is essential to require that all fields are irreducible representations of the full Poincaré group, including the discrete space and time reversal symmetries; as well as covariant under reversion, a requirement that is specific to any theory formulated with Clifford numbers.

Under these conditions, many of the difficulties encountered when using non-commutative Clifford numbers instead of real or complex numbers can be overcome. Moreover, most (if not all) apparently new results can be interpreted using standard concepts and without introducing “new” physics.

Finally, we stress again the very interesting fact that the use of the most simple Clifford algebra directly applicable to quantum physics, i.e., the biquaternions, leads in a natural way to the existence of isospin doublets, $`SU(2)`$ gauge fields, and other key features of the current Standard Model. In this respect, the predictive power of Lanczos’s equation and its generalizations by Einstein and others, as well as the many opportunities opened by the use of biquaternions and Clifford numbers in physics, are most remarquable.(22)

REFERENCES

1. W.R. Davis et al., eds., *Cornelius Lanczos Collected Published Papers With Commentaries* (North Carolina State University, Raleigh, 1998). Web page www.physics.ncsu.edu/lanczos.
2. A. Gsponer and J.-P. Hurni, “Lanczos-Einstein-Petiau: From Dirac’s equation to non-linear wave mechanics,” in Ref.1, pages 2-1248 to 2-1277.
3. C.W. Kilmister, *Proc. Roy. Irish Acad.* 57, 37–52 (1955).
4. C. Lanczos, *Zeits. f. Phys.* 57, 447–473, 474–483, 484–493 (1929). Reprinted and translated in Ref.1, pages 2-1132 to 2-1225.
5. J.D. Edmonds, Jr., *Int. J. Theor. Phys.*  6, 205–224 (1972); *Found. Phys.*  3, 313–319 (1973).
6. C. Lanczos, *Zeits. f. Phys.* 37, 405–513 (1926). Reprinted and translated in Ref.1, pages 2-932 to 2-949.
7. E.C.G. Stueckelberg et al., *Helv. Phys. Acta* 33, 727–752 (1960); 34, 621–628 (1961). See also section 2.6 in S.L. Adler, *Quaternionic Quantum Mechanics and Quantum Fields* (Oxford University Press, New York, 1995).
8. S.B.M. Bell, J.P. Cullerne, and B.M. Diaz, *Found. Phys.* 30, 35–57 (2000).
9. S. De Leo and W.A. Rodrigues, *Int. J. Theor. Phys.* 37, 1511–1529, 1707–1720 (1998).
10. J. Keller and Z. Oziewicz, eds.,  *Proc. of the Int. Conf. on the Theory of the Electron*, Adv. in Appl. Clifford Algebras 7(S) (1997).
11. S. Gull et al., *Found. Phys.* 23, 1175–1201 (1993).
12. P. Lounesto, *Found. Phys.* 23, 1203–1237 (1993).
13. F. Gürsey, *Nuovo Cim.* 7, 411–415 (1958).
14. F. Gürsey, PhD thesis, (University of London, 1950). See also F. Gürsey, *Phys. Rev.* 77, 844 (1950).
15. D. Hestenes, *Space-Time Algebra* (Gordon and Breach, New York, 1966, 1987, 1992); *J. Math. Phys.* 8, 798–808, 809–812 (1967).
16. A. Einstein and W. Mayer, *Sitzber. der Preuss. Akad*, 522–550 (1932); *Proc. Roy. Acad. Amsterdam* 36, 497–516, 615–619 (1933).
17. See, for example, T.P. Cheng and L.F. Li, *Gauge theory of elementary particle physics* (Clarendon Press, Oxford, 1984).
18. F. Gürsey, *Nuovo Cim.* 16, 230–240 (1960).
19. See, for example, J.D. Bjorken, *Phys. Rev.* D 19, 335–346 (1979).
20. G. Petiau, *Nuovo Cim.* 40, 84–101 (1965); *Ann. Inst. Henri Poincaré* 36, 89–125 (1982).
21. M. Veltman, *Reflexions on the Higgs system*, CERN Report 97-05 (1997).
22. F.J. Dyson, “Missed opportunities,” *Bull. Am. Math. Soc.* 78, 635–652 (1972); “The threefold way,” *J. Math. Phys.* 6, 1199–1215 (1962).


text:
# UWThPh-2002-4 Dissipation in a 2-dimensional Hilbert space: Various forms of complete positivityThis research was performed within the FWF Project No. P14143–PHY of the Austrian Science Foundation.

## 1 Introduction

Particle physics is not only a field where fundamental interactions are explored, but it has also become a testing ground for possible deviations from the quantum-mechanical time evolution. The time evolution of the density matrix $`\rho `$ is given by the master equation

$$\frac{d\rho }{dt}=iH\rho +i\rho H^{}D[\rho ],$$
(1)
where $`D[\rho ]`$ is a dissipative term which adds to the quantum-mechanical term on the right-hand side of Eq. (1). Such a term can emerge if the system under consideration is not fully closed but interacts weakly with the environment. In general the nature of such an interaction is unknown, but experimental data from suitable systems can be used to place bounds on the parameters of such hypothesized interactions. In the general case, the Weisskopf–Wigner approximation  allows to incorporate also unstable particles, by using non-hermitian Hamiltonians $`H=Mi\mathrm{\Gamma }/2`$, where $`M`$ and $`\mathrm{\Gamma }0`$ are hermitian.

It is not only necessary that the time evolution (1) respects $`\rho (t)0`$ $`t`$, but the stronger requirement of *complete positivity*  seems to be a natural and physical concept (for the general structure of completely positive maps, see Refs. ). This concept is defined in the following way. Let us assume that the system under consideration is described by elements of the finite complex Hilbert space $``$ ($`dimd<\mathrm{}`$) with time evolution $`\rho (t)\gamma _t(\rho )`$ with $`\rho (0)=\rho `$. Considering in addition the finite-dimensional Hilbert space $`^n`$, one can extend the time evolution $`\gamma _t`$ on $``$ to a time evolution $`\gamma _{n;t}`$ on $`^n`$ by defining $`\gamma _{n;t}=\gamma _t\mathrm{𝟙}`$. If the time evolution $`\gamma _t`$ derived from Eq. (1) with the dissipative term $`D[\rho ]`$ has the property that $`\gamma _{n;t}(\rho )0`$ is valid for all times $`t`$, all $`n=0,1,2,\mathrm{}`$ and all density matrices on the space $`^n`$, then $`\gamma _t`$ is called completely positive.

Complete positivity of $`\gamma _t`$ determines the general structure of $`D[\rho ]`$  (see also Refs. ). It has been shown by Lindblad  (see also Ref. ) that $`\gamma _t`$ of Eq. (1) is completely positive if and only if $`D[\rho ]`$ has the structure

$`D[\rho ]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{r}{}}}\left(A_j^{}A_j\rho +\rho A_j^{}A_j2A_j\rho A_j^{}\right)`$ (2)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=1}{\overset{r}{}}}\left([A_j^{},A_j\rho ]+[\rho A_j^{},A_j]\right),`$
where the operators $`A_j`$ act on $``$.

Using the relation

$$A_j=B_j+s_j\mathrm{𝟙}\mathrm{with}\mathrm{Tr}B_j=0,$$
(3)
the dissipative term (2) can be reformulated as

$$D[\rho ]=i[\mathrm{\Delta }H,\rho ]+D^{}[\rho ]\mathrm{with}\mathrm{\Delta }H=\frac{i}{2}\underset{j=1}{\overset{r}{}}\left(s_jB_j^{}s_j^{}B_j\right),$$
(4)
where $`D^{}[\rho ]`$ is obtained from $`D[\rho ]`$ by the replacement $`A_jB_j`$. This reformulation has the effect that part of $`D[\rho ]`$ is shifted into the quantum-mechanical term of the time evolution (1), such that a new Hamiltonian $`H^{}=H+\mathrm{\Delta }H`$ appears. Note that for hermitian operators $`A_j`$ we have $`\mathrm{\Delta }H=0`$. In the space of traceless operators on the Hilbert space $``$ we can choose a basis $`\left\{F_j\right|j=1,\mathrm{},d^21\}`$ with the property

$$\mathrm{Tr}\left(F_j^{}F_k\right)=\delta _{jk}$$
(5)
and expand the operators $`B_j`$ as

$$B_j=\underset{k=1}{\overset{d^21}{}}C_{kj}F_k.$$
(6)
Then we obtain the expression

$$D^{}[\rho ]=\frac{1}{2}\underset{j=1}{\overset{r}{}}\left([B_j,\rho B_j^{}]+[B_j\rho ,B_j^{}]\right)=\frac{1}{2}\underset{k,l=1}{\overset{d^21}{}}c_{kl}\left([F_k,\rho F_l^{}]+[F_k\rho ,F_l^{}]\right),$$
(7)
where $`(c_{kl})`$ is a positive matrix defined by

$$c_{kl}=\underset{j=1}{\overset{r}{}}C_{kj}C_{lj}^{}.$$
(8)
The form (7) of the dissipative term has been derived by Gorini, Kossakowski and Sudarshan . It is equivalent to the Lindblad form (2). Thus the time evolution with $`H`$ and $`D[\rho ]`$ is equivalent to the one with $`H^{}`$ and $`D^{}[\rho ]`$.

In particle physics, searches for deviations from the quantum-mechanical time evolution are going on in neutral meson-antimeson systems ($`K^0\overline{K}^0`$ and $`B^0\overline{B}^0`$) (for a list of papers see, e.g., Refs. ) and neutrino physics (for a list of papers see, e.g., Refs. ). The importance of complete positivity, in particular, in $`K^0\overline{K}^0`$ and analogous systems, has been stressed in Refs. . For instance, only if $`\gamma _t`$ is completely positive, the positivity of the time evolution $`\gamma _t\gamma _t`$ in the space $``$ is guaranteed. This follows from the decomposition $`\gamma _t\gamma _t=(\gamma _t\mathrm{𝟙})(\mathrm{𝟙}\gamma _t)`$, where both factors are positive according to complete positivity. Forming the tensor product $`\gamma _t\gamma _t`$ is a method for implementing the 1-particle time evolution at the 2-particle level, which is an often used procedure, e.g., in the $`K^0\overline{K}^0`$ system. For an example where $`\gamma _t`$ is only positive but not completely positive, with ensuing non-positivity of $`\gamma _t\gamma _t`$, see Ref. .

For simplicity we assume a hermitian Hamiltonian $`H`$ from now on, but this assumption is irrelevant for our discussion of complete positivity; it only concerns the investigation of the asymptotic limit of the time evolution, where loss of probability due to a non-hermitian $`H`$ as obtained in the Weisskopf–Wigner approximation  will lead to $`\rho (t)0`$ for $`t\mathrm{}`$.

Apart from complete positivity of the time evolution, the assumptions on which the present work is based are the following:

1. We work in a 2-dimensional complex Hilbert space, which we can identify with $`^2`$;
2. we assume hermitian Lindblad operators, i.e., $`A_j=A_j^{}`$$`j`$.

Both assumptions are crucial for the following discussions. The first one is motivated by the applications in particle physics, whereas the second assumption guarantees that the entropy $`S[\rho ]=\text{Tr}(\rho \mathrm{ln}\rho )`$ cannot decrease as a function of time . In this framework we will discuss the different forms of the dissipative term $`D[\rho ]`$ used in the literature and we will show their equivalence. We will put emphasis on the formulation of the time evolution (1) in $`^3`$, where we will represent $`D[\rho ]`$ as a positive sum over projectors onto planes. We will also study in detail the matrix formulation of $`D[\rho ]`$ as advocated by Benatti and Floreanini, e.g., in Refs. , and relate this formulation with the geometric version of $`D[\rho ]`$ as a sum of projectors. Finally, we will investigate the limit $`t\mathrm{}`$ of the density matrix.

## 2 Equivalent forms of the dissipative term

We start with the original form of the dissipative term (2) and take into account that we confine ourselves to hermitian Lindblad operators $`A_j`$. Thus $`D[\rho ]`$ simplifies to

$$\text{Form A:}D[\rho ]=\frac{1}{2}\underset{j=1}{\overset{r}{}}\left(A_j^2\rho +\rho A_j^22A_j\rho A_j\right)=\frac{1}{2}\underset{j=1}{\overset{r}{}}[A_j,[A_j,\rho ]].$$
(9)
The number of Lindblad terms in Eq. (9) is denoted by $`r`$.

Next we note that the dissipative term can be rewritten in terms of projectors $`P_j`$ and their orthogonal complements $`P_j^{}=\mathrm{𝟙}P_j`$ as

$$\text{Form B:}D[\rho ]=\frac{1}{2}\underset{j=1}{\overset{r}{}}\lambda _j\left(P_j\rho P_j^{}+P_j^{}\rho P_j\right).$$
(10)
The projectors $`P_j`$ are non-trivial projectors in $`^2`$, which can be parameterized as

$$P_j=\frac{1}{2}\left(\mathrm{𝟙}+\stackrel{}{n}_j\stackrel{}{\sigma }\right),$$
(11)
where the $`\stackrel{}{n}_j`$ are real unit vectors and $`\stackrel{}{\sigma }`$ denotes the vector of Pauli matrices. The quantities $`\lambda _j`$ are real, positive numbers.

###### Proposition 1

Forms A and B of $`D[\rho ]`$ are equivalent.

Proof: A general hermitian $`2\times 2`$ matrix $`A_j`$ can be represented by

$$A_j=\frac{1}{2}\left(a_j\mathrm{𝟙}+\sqrt{\lambda _j}\stackrel{}{n}_j\stackrel{}{\sigma }\right),$$
(12)
where $`a_j`$ and $`\lambda _j`$ are real numbers ($`\lambda _j0`$) and $`\stackrel{}{n}_j`$ is a unit vector. Note that the part of $`A_j`$ proportional to the unit matrix does not contribute to $`D[\rho ]`$, as evident from Eq. (9). Therefore, $`\lambda _j`$ must not be zero in order to have a non-trivial effect of $`A_j`$. Consequently, we are allowed to replace $`A_j`$ by $`\sqrt{\lambda _j}P_j`$ in $`D[\rho ]`$. Using Eq. (9) and $`P_j^2=P_j`$, we derive

$$D[\rho ]=\frac{1}{2}\underset{j=1}{\overset{r}{}}\lambda _j\left(P_j\rho +\rho P_j2P_j\rho P_j\right),$$
(13)
which can be rewritten in Form B of Eq. (10). $`\mathrm{}`$

The time evolution (1) is easily reformulated as a differential equation for a real 3-vector $`\stackrel{}{\rho }`$ by using

$$\rho =\frac{1}{2}\left(\mathrm{𝟙}+\stackrel{}{\rho }\stackrel{}{\sigma }\right)\mathrm{and}H=\frac{1}{2}\left(\mathrm{𝟙}+\stackrel{}{h}\stackrel{}{\sigma }\right).$$
(14)
Since for simplicity we have assumed that $`H`$ is hermitian, the 3-vector $`\stackrel{}{h}`$ is real as well (for an extension to non-hermitian Hamiltonians see, e.g., Ref. ). The new version of Eq. (1) is given by

$$\frac{\mathrm{d}\stackrel{}{\rho }}{\mathrm{d}t}=\stackrel{}{h}\times \stackrel{}{\rho }L\stackrel{}{\rho }.$$
(15)
With Form A or B of the dissipative term and Eq. (12) or (11), we arrive at a geometric version of $`D[\rho ]`$

$$\text{Form C:}L\stackrel{}{\rho }=\frac{1}{2}\underset{j=1}{\overset{r}{}}\lambda _j\left(\stackrel{}{\rho }\stackrel{}{n}_j\stackrel{}{n}_j\stackrel{}{\rho }\right),$$
(16)
where the matrix $`L`$ is a positive linear combination of projectors

$$𝒫(\stackrel{}{n}_j)=\mathrm{𝟙}_3\stackrel{}{n}_j\stackrel{}{n}_j^T.$$
(17)
The projector (17) projects onto the plane orthogonal to $`\stackrel{}{n}_j`$. Clearly, Form C is equivalent to the the previous forms of the dissipative term.

The dissipation matrix $`L`$ in Form C can be reformulated as

$$\text{Form D:}L_{\alpha \beta }=\frac{1}{2}\left(\mathrm{\Lambda }\delta _{\alpha \beta }𝐪_\alpha 𝐪_\beta \right)$$
(18)
with

$$𝐪_\alpha ^r(\alpha =1,2,3)\mathrm{and}\mathrm{\Lambda }=\underset{\alpha =1}{\overset{3}{}}|𝐪_\alpha |^2.$$
(19)
###### Proposition 2

Forms C and D are equivalent.

Proof: The proof of this statement amounts to a mere rewriting of the elements of $`L`$ by

$$\underset{j=1}{\overset{r}{}}\lambda _j\left(\delta _{\alpha \beta }(\stackrel{}{n}_j)_\alpha (\stackrel{}{n}_j)_\beta \right)=\underset{j=1}{\overset{r}{}}\lambda _j\delta _{\alpha \beta }\underset{j=1}{\overset{r}{}}\sqrt{\lambda _j}(\stackrel{}{n}_j)_\alpha \sqrt{\lambda _j}(\stackrel{}{n}_j)_\beta .$$
(20)
Then we define

$$𝐪_\alpha =\left(\begin{array}{c}\sqrt{\lambda _1}(\stackrel{}{n}_1)_\alpha \\ \mathrm{}\\ \sqrt{\lambda _r}(\stackrel{}{n}_r)_\alpha \end{array}\right)\mathrm{and}\mathrm{\Lambda }=\underset{j=1}{\overset{r}{}}\lambda _j.$$
(21)
The sum over the square of the lengths of the three vectors $`𝐪_\alpha `$ is performed via

$$\underset{\alpha =1}{\overset{3}{}}|𝐪_\alpha |^2=\underset{\alpha =1}{\overset{3}{}}\underset{j=1}{\overset{r}{}}\lambda _j(\stackrel{}{n}_j)_\alpha ^2=\underset{j=1}{\overset{r}{}}\lambda _j\left(\underset{\alpha =1}{\overset{3}{}}(\stackrel{}{n}_j)_\alpha ^2\right)=\underset{j=1}{\overset{r}{}}\lambda _j.$$
(22)
Here, we have used that the $`\stackrel{}{n}_j`$ are 3-dimensional unit vectors. Thus we have obtained Form D of $`L`$ from Form C. This procedure can be reversed: given three vectors $`𝐪_\alpha ^r`$, we can use Eq. (21) to construct $`r`$ unit vectors $`\stackrel{}{n}_j`$ and positive numbers $`\lambda _j`$. In this way, we gain Form C from Form D. $`\mathrm{}`$

The question arises how many terms are necessary in $`D[\rho ]`$ in the most general case. This question is answered by the following theorem.

###### Theorem 1

The most general case is covered by three Lindblad terms in $`D[\rho ]`$, i.e., if $`D[\rho ]`$ is given by a sum over more than three terms, then it can be rewritten as a sum of at most three terms. If $`D[\rho ]`$ is formulated with the minimal number of terms, then, using Forms B or C, there are three distinct minimal cases referring to one, two or three linearly independent vectors $`\stackrel{}{n}_j`$ with $`\lambda _j>0`$; we will denote these cases by an *index* $`\mathrm{}\{1,2,3\}`$.Obviously, in the minimal formulation of $`D[\rho ]`$ we have $`r\mathrm{}`$.

Proof: Let us first assume that $`r>3`$. Since there are only three vectors $`𝐪_\alpha `$, we can find a rotation $`R`$ acting on $`^r`$ such that

$$R𝐪_\alpha =\left(\begin{array}{c}𝐐_\alpha \\ 0\\ \mathrm{}\\ 0\end{array}\right)\mathrm{with}𝐐_\alpha ^3.$$
(23)
In the most general case the set of vectors $`Q=\{𝐐_\alpha |\alpha =1,2,3\}`$ is linearly independent and can be parameterized by

$$𝐐_\alpha =\left(\begin{array}{c}q_\alpha ^1\\ q_\alpha ^2\\ q_\alpha ^3\end{array}\right)=\left(\begin{array}{c}\sqrt{\mu _1}(\stackrel{}{m}_1)_\alpha \\ \sqrt{\mu _2}(\stackrel{}{m}_2)_\alpha \\ \sqrt{\mu _3}(\stackrel{}{m}_3)_\alpha \end{array}\right).$$
(24)
The real and positive numbers $`\mu _j`$ are chosen in such a way that the vector $`\stackrel{}{m}_1`$, extracted from the first elements of the vectors $`𝐐_\alpha `$, is a unit vector; the same is done for the second and third elements.§§§This procedure is analogous to the one to prove that Form C follows from Form D. In this case we have an index $`\mathrm{}=3`$. If $`Q`$ spans a 2-dimensional space, we choose the rotation $`R`$ such that the third elements of all $`𝐐_\alpha `$ are zero; if $`Q`$ spans a 1-dimensional space, the second and third elements are taken to be zero. These two cases refer to index $`\mathrm{}=2`$ and 1, respectively. Form D of the dissipative term tells us that $`D[\rho ]`$ is independent of any rotation $`R`$. Consequently, we arrive at

$$\underset{j=1}{\overset{r}{}}\lambda _j\left(\delta _{\alpha \beta }(\stackrel{}{n}_j)_\alpha (\stackrel{}{n}_j)_\beta \right)=\underset{j=1}{\overset{\mathrm{}}{}}\mu _j\left(\delta _{\alpha \beta }(\stackrel{}{m}_j)_\alpha (\stackrel{}{m}_j)_\beta \right),$$
(25)
which proves the theorem for $`r>3`$. For $`r3`$ we follow analogous steps performed after Eq. (23), with the rotation $`R`$ acting now on a 3 or 2-dimensional space. For $`r=1`$ the procedure is trivial. $`\mathrm{}`$

On the other hand, to find the minimal number of operators $`A_j`$ needed for the most general dissipative term (2), we could start from the Gorini-Kossakowski-Sudarshan expression (7). Transforming back to the Lindblad structure (2) by using the relations

$$c_{kl}=\left(U\widehat{c}U^{}\right)_{kl}\mathrm{and}B_j=\underset{k=1}{\overset{d^21}{}}U_{kj}\sqrt{\widehat{c}_{jj}}F_k,$$
(26)
where $`\widehat{c}0`$ is diagonal and $`U`$ a unitary matrix, it can be seen that, in the general case, we obtain $`d^21`$ terms in $`D[\rho ]`$ of Eq. (2); if some of the elements $`\widehat{c}_{jj}`$ are zero, we will have less than $`d^21`$ terms. Applying this to $`d=2`$, we find at most three terms, which agrees with the result of the explicit calculations leading to Theorem 1.

## 3 Complete-positivity conditions on the dissipation matrix $`L`$

Benatti and Floreanini  parameterized the dissipation matrix $`L`$ by 6 real constants and expressed complete positivity in the form of inequalities satisfied by these parameters. Thus they have the version

$$\text{Form E:}L=2\left(\begin{array}{ccc}a& b& c\\ b& \alpha & \beta \\ c& \beta & \gamma \end{array}\right),$$
(27)
together with

$`\begin{array}{c}2R\alpha +\gamma a0,\\ 2Sa+\gamma \alpha 0,\\ 2Ta+\alpha \gamma 0;\end{array}`$ (28d)
$`RSb^2,RTc^2,ST\beta ^2;`$ (28e)
$`RST2bc\beta +R\beta ^2+Sc^2+Tb^2.`$ (28f)
We will show now that Form E is equivalent to the forms presented in the previous section.

First we want to put Eqs. (28d), (28e), (28f) into a simpler equivalent form.

###### Lemma 1

Given $`L`$, there exists a symmetric matrix $`M`$ such that

$$L=\frac{1}{2}(\mathrm{Tr}M\mathbb{\hspace{0.17em}1}_3M).$$
(29)
In terms of $`M`$, Eqs. (28d), (28e), (28f) are given by

$$\begin{array}{ccc}\hfill \text{(i)}& M_{\alpha \alpha }0& (\alpha =1,2,3),\hfill \\ \hfill \text{(ii)}& M_{\alpha \alpha }M_{\beta \beta }M_{\alpha \beta }^2& \alpha \beta ,\hfill \\ \hfill \text{(iii)}& detM0,& \end{array}$$
(30)
respectively.

Proof: For $`\alpha \beta `$ we have $`M_{\alpha \beta }=2L_{\alpha \beta }`$. The diagonal elements of $`M`$ are determined by $`M_{\alpha \alpha }=L_{\alpha \alpha }+L_{\beta \beta }+L_{\gamma \gamma }`$, where $`\alpha \beta \gamma \alpha `$. The second part of the lemma is proven by plugging Eq. (29) into Eqs. (28d), (28e), (28f). $`\mathrm{}`$

The next lemma will allow us to connect $`M`$ possessing properties (30) with Form D of the dissipative term.

###### Lemma 2

A real and symmetric $`3\times 3`$ matrix $`M`$ has the properties of Eq. (30) if and only if there exist vectors $`𝐪_\alpha ^3`$ ($`\alpha =1,2,3`$) such that

$$M_{\alpha \beta }=𝐪_\alpha 𝐪_\beta .$$
(31)
Proof:
($``$) This direction of the proof is quickly dealt with. Since $`M_{\alpha \alpha }=𝐪_\alpha ^20`$, property (i) is valid. Furthermore, for $`\alpha \beta `$, with the Cauchy–Schwarz inequality we derive $`M_{\alpha \beta }^2=(𝐪_\alpha 𝐪_\beta )^2𝐪_\alpha ^2𝐪_\beta ^2=M_{\alpha \alpha }M_{\beta \beta }`$ and property (ii) holds as well. Defining a $`3\times 3`$ matrix $`q=(𝐪_1,𝐪_2,𝐪_3)`$, we have $`detM=det(q^Tq)=(detq)^20`$ and thus property (iii). This completes the first half of the proof.
($``$) This direction of the proof is more involved. If all $`M_{\alpha \alpha }`$ were zero, then according to property (ii) of Eq. (30) we would have $`M=0`$. Therefore, at least one of the diagonal elements of $`M`$ is non-zero. Without loss of generality we assume $`M_{11}>0`$. First we consider the case

$$M_{11}M_{22}M_{12}^2>0.$$
(32)
Denoting the elements of $`𝐪_\alpha `$ by $`q_\alpha ^j`$, we can define a vector $`𝐪_1`$ by $`q_1^1=\sqrt{M_{11}}`$, $`q_1^2=q_1^3=0`$. From $`M_{12}=𝐪_1𝐪_2`$, it follows after a sign choice that $`q_2^1=M_{12}/\sqrt{M_{11}}`$. Taking into account that $`𝐪_2𝐪_2=M_{22}`$ and defining $`q_2^3=0`$, the first two vectors are given by

$$𝐪_1=\left(\begin{array}{c}\sqrt{M_{11}}\\ 0\\ 0\end{array}\right),𝐪_2=\left(\begin{array}{c}M_{12}/\sqrt{M_{11}}\\ \sqrt{M_{22}M_{12}^2/M_{11}}\\ 0\end{array}\right).$$
(33)
With the relation of Eq. (32) we find that $`𝐪_2`$ is a well-defined real 3-vector. Next we use $`M_{13}=𝐪_1𝐪_3`$ and $`M_{23}=𝐪_2𝐪_3`$ and obtain

$$q_3^1=\frac{M_{13}}{\sqrt{M_{11}}}\text{and}q_3^2=\frac{M_{23}M_{12}M_{13}/M_{11}}{\sqrt{M_{22}M_{12}^2/M_{11}}}.$$
(34)
It remains to take into account $`𝐪_3𝐪_3=M_{33}`$. After some algebra we arrive at

$$(q_3^3)^2=M_{33}(q_3^1)^2(q_3^2)^2=\frac{detM}{M_{11}M_{22}M_{12}^2}0.$$
(35)
The positivity follows from property (iii) of Eq. (30). Thus, $`q_3^3`$ is well-defined and we have proven relation (31), provided condition (32) holds.

It remains to check the same for the special case

$$M_{11}M_{22}=M_{12}^2,$$
(36)
which was excluded by Eq. (32). From Eq. (36) we obtain $`M_{12}=\eta \sqrt{M_{11}M_{22}}`$ with $`\eta =\pm 1`$ and

$$detM=\left(\sqrt{M_{11}}M_{23}\eta \sqrt{M_{22}}M_{13}\right)^2.$$
(37)
Since $`detM0`$ (see Eq. (30)), it follows that

$$M_{23}=\eta M_{13}\sqrt{M_{22}/M_{11}}.$$
(38)
With this relation and taking into account condition (36), it is easy to check that

$$𝐪_1=\sqrt{M_{11}}\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right),𝐪_2=\eta \sqrt{M_{22}}\left(\begin{array}{c}1\\ 0\\ 0\end{array}\right),𝐪_3=\left(\begin{array}{c}M_{13}/\sqrt{M_{11}}\\ \sqrt{M_{33}M_{13}^2/M_{11}}\\ 0\end{array}\right)$$
(39)
represents a consistent choice of vectors, which fulfills equation (31). This completes the proof. $`\mathrm{}`$

With Lemmata 1 and 2 we now readily see that a matrix $`L`$ fulfills the conditions of Eqs. (28d), (28e), (28f) if and only if $`L`$ is given by $`L_{\alpha \beta }=\frac{1}{2}(\mathrm{\Lambda }\delta _{\alpha \beta }𝐪_\alpha 𝐪_\beta )`$, where $`𝐪_\alpha ^3`$ and $`\mathrm{\Lambda }=_{\alpha =1}^3|𝐪_\alpha |^2`$.

###### Proposition 3

Forms D and E of the dissipation matrix $`L`$ are equivalent.

Proof: According to Lemma 1, from the matrix $`L`$ we construct a matrix $`M`$ with properties (30). Lemma 2 tells us that such an $`M`$ is represented by a matrix of scalar products (see Eq. (31)) and vice versa. Therefore, our statement is true. $`\mathrm{}`$

## 4 The asymptotic limit $`t\mathrm{}`$

Now we consider the asymptotic limit of the density matrix $`\rho (t)`$ with time evolution (1) . For this purpose we use the results of Theorem 1, where we have also defined the index $`\mathrm{}`$ of $`L`$. For a general study of the large time behaviour of the density matrix starting with expression (7) see Ref. .

###### Theorem 2

For index $`\mathrm{}=2`$ or $`\mathrm{}=3`$, the asymptotic limit of the density matrix is given by

$$\underset{t\mathrm{}}{lim}\rho (t)=\frac{1}{2}\mathrm{𝟙}.$$
(40)
Proof: This statement is most easily proven by using Form C, Eq. (16), of the dissipative term and Eq. (15) of the time evolution. Denoting by $`𝒜`$ the operator on the right-hand side of Eq. (15), we have

$$𝒜𝐱=\stackrel{}{h}\times 𝐱\frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}\lambda _j𝒫(\stackrel{}{n}_j)𝐱$$
(41)
for arbitrary complex 3-vectors $`𝐱`$. If we can show that every eigenvalue $`c`$ of $`𝒜`$ fulfills $`\mathrm{Re}c<0`$, the theorem is proven. Let $`𝐱`$ be a normalized eigenvector with eigenvalue $`c`$. Then, with $`𝐱^{}𝐱=1`$ we have

$$c=𝐱^{}𝒜𝐱=2i\stackrel{}{h}(\mathrm{Re}𝐱\times \mathrm{Im}𝐱)\frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}\lambda _j𝐱^{}𝒫(\stackrel{}{n}_j)𝐱,$$
(42)
and, consequently,

$$\mathrm{Re}c=\frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}\lambda _j𝐱^{}𝒫(\stackrel{}{n}_j)𝐱.$$
(43)
Since the projectors $`𝒫(\stackrel{}{n}_j)`$ are positive operators and $`\lambda _j>0`$, we have $`\mathrm{Re}c0`$. If $`\mathrm{Re}c=0`$, it is necessary that $`𝐱^{}𝒫(\stackrel{}{n}_j)𝐱=1|\stackrel{}{n}_j𝐱|^2=0`$ $`j=1,\mathrm{},\mathrm{}`$. Therefore, the eigenvector $`𝐱`$ is proportional to $`\stackrel{}{n}_j`$ $`j=1,\mathrm{},\mathrm{}`$. This is a contradiction for $`\mathrm{}=2`$ or 3 independent vectors $`\stackrel{}{n}_j`$ and, indeed, every eigenvalue $`c`$ has a negative non-zero real part. $`\mathrm{}`$

###### Theorem 3

For $`\mathrm{}=1`$ ($`\stackrel{}{n}_1\stackrel{}{n}`$, $`P_1P`$, $`\lambda _1\lambda `$) we have either $`\stackrel{}{h}\stackrel{}{n}`$ or, equivalently, $`[H,P]=0`$, in which case we obtain

$$\underset{t\mathrm{}}{lim}\rho (t)=P\rho (0)P+P^{}\rho (0)P^{};$$
(44)
or $`\stackrel{}{h}\parallel ̸\stackrel{}{n}`$, i.e., $`[H,P]0`$, then the asymptotic limit of $`\rho (t)`$ is the same as in Theorem 2.

Proof: We follow the same strategy as in the proof of the previous theorem. Thus either $`\mathrm{Re}c<0`$ or the eigenvector $`𝐱`$ is proportional to $`\stackrel{}{n}`$, in which case $`\mathrm{Re}c=0`$. Assuming $`𝐱=\stackrel{}{n}`$ without loss of generality, we have now $`𝒜𝐱=\stackrel{}{h}\times \stackrel{}{n}=c\stackrel{}{n}`$, where the eigenvalue $`c`$ is imaginary. This equation is only soluble for $`\stackrel{}{h}\stackrel{}{n}`$, whence it follows that $`c=0`$. The relation $`\stackrel{}{h}\stackrel{}{n}`$ is equivalent to $`[H,P]=0`$. In this case we can write the Hamiltonian as $`H=hP+h^{}P^{}`$. Decomposing an arbitrary density matrix $`\rho `$ as

$$\rho =\rho _0+\rho _1\mathrm{with}\rho _0=P\rho P+P^{}\rho P^{}\mathrm{and}\rho _1=P\rho P^{}+P^{}\rho P,$$
(45)
we find

$$\frac{d\rho _0}{dt}=0\text{and}\frac{d\rho _1}{dt}=[i(hh^{})\lambda /2]P\rho P^{}+[i(hh^{})\lambda /2]P^{}\rho P.$$
(46)
Therefore, we arrive at

$$\rho _0(t)=\rho _0(0)\text{and}\rho _1(t)=e^{[i(hh^{})\lambda /2]t}P\rho (0)P^{}+e^{[i(hh^{})\lambda /2]t}P^{}\rho (0)P.$$
(47)
Since $`lim_t\mathrm{}\rho _1(t)=0`$, the theorem is proven. $`\mathrm{}`$

The different limits of $`\rho (t)`$ discussed in Theorems 2 and 3 have been noticed in Refs. . In the case of $`\mathrm{}=1`$ and $`[H,P]=0`$, the limit (44) of the density matrix has the form $`\rho =\mu P+(1\mu )P^{}`$ with $`0\mu 1`$; all density matrices obeying $`d\rho /dt=0`$ have this form for $`[H,P]=0`$. For $`\mathrm{}=1`$ and $`[H,P]0`$, and for $`\mathrm{}=2,3`$, the unique density matrix which is time-independent is proportional to the unit matrix.

## 5 Summary and Conclusions

In this paper we have considered the quantum-mechanical time evolution of a $`2\times 2`$ density matrix $`\rho (t)`$, where the von Neumann equation is modified by a dissipative term $`D[\rho ]`$, which, therefore, must be of the Lindblad type or, equivalently, of the Gorini-Kossakowski-Sudarshan type in order to assure a completely positive time evolution. We have, furthermore, assumed that the Lindblad operators $`A_j`$ are all hermitian which ensues that the entropy is non-decreasing with time. Our starting point conforms with many applications of the open-systems approach in particle physics.

We have discussed five equivalent forms of $`D[\rho ]`$ which all have their merits depending on the problem considered. We have put particular emphasis on the time evolution in the form of Eq. (15), where the density matrix and the Hamiltonian are represented by real 3-vectors, and Form C, Eq. (16), of the dissipative term, where the dissipative term is a positive linear combination of projectors onto 2-dimensional planes in $`^3`$.

We have studied the question of the minimal number $`\mathrm{}`$ of Lindblad terms needed in order to reproduce a given $`D[\rho ]`$ and formulated the result in Theorem 1; the proof of this theorem represents at the same time a procedure how to determine $`\mathrm{}`$ in practice. An other procedure would be to start with the Gorini-Kossakowski-Sudarshan expression from which it also can be seen that $`D[\rho ]`$ can be generally decomposed into 3 terms.

We have also connected the approach where the dissipative term is given by a matrix $`L`$ specified by the conditions (28d), (28e), (28f) with the geometric picture of $`D[\rho ]`$ as given by Form C, Eq. (16). Again, the proof which shows the equivalence between the two approaches, given by Lemmata 1 and 2, indicates a practical way to obtain the projectors $`𝒫(\stackrel{}{n}_j)`$ (17) associated with $`D[\rho ]`$.

Finally, we have presented a general discussion of the limit $`t\mathrm{}`$ of $`\rho (t)`$, where the usefulness of Form C of $`D[\rho ]`$ was exemplified.