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# Asymptotic expansions close to the singularity in Gowdy spacetimes

## 1. Introduction

The motivation for studying the problem discussed in this article is the desire to understand the structure of singularities in cosmological spacetimes. By the singularity theorems, cosmological spacetimes typically have a singularity in the sense of causal geodesic incompleteness. However, it seems that the methods used to obtain this result are not so well suited to answering related questions concerning e.g. curvature blow up. To proceed, it seems difficult to avoid analyzing the equations in detail. After some appropriate choice of gauge, one is then confronted with the task of analyzing the asymptotics of a non-linear hyperbolic equation. This is in general not so easy. Consequently, one often imposes some symmetry condition, and we will here consider a class of spacetimes with a two dimensional group of symmetries. The problem one ends up with is then a system of non-linear wave equations in $`1+1`$ dimensions.

The Gowdy spacetimes were first introduced in  (see also ), and in  the fundamental questions concerning global existence were answered. We will take the Gowdy vacuum spacetimes on $`\times T^3`$ to be metrics of the form (1). We do not wish to motivate this choice at length here, but refer the reader to the above mentioned references for further details. A brief justification is given in the introduction of . Let

(1)
$$g=e^{(\tau \lambda )/2}(e^{2\tau }d\tau ^2+d\theta ^2)+e^\tau [e^Pd\sigma ^2+2e^PQd\sigma d\delta +(e^PQ^2+e^P)d\delta ^2]$$

Here, $`\tau `$ and $`(\theta ,\sigma ,\delta )`$ are coordinates on $`T^3`$. The evolution equations become

(2) $`P_{\tau \tau }e^{2\tau }P_{\theta \theta }e^{2P}(Q_\tau ^2e^{2\tau }Q_\theta ^2)`$ $`=`$ $`0`$
(3) $`Q_{\tau \tau }e^{2\tau }Q_{\theta \theta }+2(P_\tau Q_\tau e^{2\tau }P_\theta Q_\theta )`$ $`=`$ $`0,`$
and the constraints

(4) $`\lambda _\tau `$ $`=`$ $`P_\tau ^2+e^{2\tau }P_\theta ^2+e^{2P}(Q_\tau ^2+e^{2\tau }Q_\theta ^2)`$
(5) $`\lambda _\theta `$ $`=`$ $`2(P_\theta P_\tau +e^{2P}Q_\theta Q_\tau ).`$
Obviously, the constraints are decoupled from the evolution equations, excepting the condition on $`P`$ and $`Q`$ implied by (5). Thus the equations of interest are the two non-linear coupled wave equations (2)-(3). In the above parametrization, the singularity corresponds to $`\tau \mathrm{}`$, and the subject of this article is the asymptotics of solutions to (2)-(3) as $`\tau \mathrm{}`$. The equations (2)-(3) are wave map equations. In fact, let

$$g_0=e^{2\tau }d\tau ^2+d\theta ^2+e^{2\tau }d\chi ^2$$

be a metric on $`\times T^2`$ and let

(6)
$$g_R=dP^2+e^{2P}dQ^2$$

be a metric on $`^2`$. Then (2)-(3) are the wave map equations for a map from $`(\times T^2,g_0)`$ to $`(^2,g_R)`$ which is independent of the $`\chi `$-coordinate. Note that $`(^2,g_R)`$ is isometric to the upper half plane $`H=\{(x,y)^2:y>0\}`$ with metric

(7)
$$g_H=\frac{dx^2+dy^2}{y^2}$$

under the map

(8)
$$\varphi _{RH}(Q,P)=(Q,e^P).$$

Thus the target space is hyperbolic space. For the rest of the article, we will be concerned with the above mentioned wave map equations and not consider the consequences for the resulting spacetimes. The implications for the spacetime geometry can be found elsewhere, see e.g. .

The idea of finding expansions for solutions close to the singularity started with the article  by Grubišić and Moncrief. In our setting, the natural expansions are

(9) $`P(\tau ,\theta )`$ $`=`$ $`v(\theta )\tau +\varphi (\theta )+e^{ϵ\tau }u(\theta ,\tau )`$
(10) $`Q(\tau ,\theta )`$ $`=`$ $`q(\theta )+e^{2v(\theta )\tau }[\psi (\theta )+w(\tau ,\theta )]`$
where $`ϵ>0`$ and $`w,u0`$ as $`\tau \mathrm{}`$ and $`0<v(\theta )<1`$. This should be compared with (5) and (6) of , where $`Z=P`$, $`X=Q`$ and $`t=e^\tau `$. A heuristic argument motivating the condition on the velocity can be found in . In the non-generic case $`Q=0`$, one can prove that (9) holds without the condition on $`v`$. This special case is called polarized Gowdy and has been studied in , which also considers the other topologies for Gowdy spacetimes. The numerical simulations indicate that there are exceptions to (9)-(10), called spikes, at which the kinetic energy of the wave map makes a jump. We refer the reader to  for numerical and to  for analytical results concerning spikes.

By the so called Fuchsian techniques one can construct a large family of solutions with the asymptotic behaviour (9)-(10). In fact, given functions $`v,\varphi ,q`$ and $`\psi `$ from $`S^1`$ to $``$ of a suitable degree of smoothness and subject to the condition $`0<v<1`$, one can construct solutions to (2)-(3) with asymptotics of the form (9)-(10). The proof of this in the real analytic case can be found in  and  covers the smooth case. One nice feature of this construction is the fact that one gets to specify four functions freely, just as as if though one were specifying initial data for (2)-(3).

Finally, in  a condition on the initial data leading to the desired asymptotics is given. Note also the related work in . In , conditions on the $`H^2`$ norm of the first derivatives are imposed. The reason is the following. It turns out that the most important quantity to control is the velocity $`v`$ appearing in (9). In other words, one wants to control what $`P_\tau `$ converges to. In , this control was achieved by integrating (2). In order to reach any conclusion, one then has to control $`e^{2\tau }P_{\theta \theta }`$ in the sup norm, which is the reason for the condition on the $`H^2`$ norm of the first derivatives. Looking at the expansion (9) naively, another possibility presents itself. One expects $`P_\tau `$ to converge to $`v`$ rapidly, so that $`P_\tau P/\tau `$ should be of the order of magnitude $`O(\tau ^1)`$. Assume for a moment that we can prove that we have this estimate in the sup norm. Then one easily sees that $`_\tau (P/\tau )=O(\tau ^2)`$. Consequently, $`P/\tau `$ converges to a continuous function. Since $`P_\tau P/\tau `$ converges to zero, we then get the convergence of $`P_\tau `$ to a continuous function. What we need is consequently to prove that $`P_\tau P/\tau `$ tends to zero as $`\tau ^1`$ under suitable circumstances. One quantity suited to achieve this task is

(11)
$$F(\tau )=\frac{1}{2}\underset{\theta S^1}{sup}[(P_\tau \frac{1}{\tau }P+e^\tau P_\theta )^2+e^{2P}(Q_\tau +e^\tau Q_\theta )^2]$$

$$+\frac{1}{2}\underset{\theta S^1}{sup}[(P_\tau \frac{1}{\tau }Pe^\tau P_\theta )^2+e^{2P}(Q_\tau e^\tau Q_\theta )^2].$$

This object may seem unnatural at first, however it appears naturally in the arguments. It turns out that if $`1P\tau 1`$ in an interval $`[\tau _1,\tau _2]`$, then $`F(\tau )F(\tau _1)(\tau _1/\tau )^2`$ for $`\tau [\tau _1,\tau _2]`$, see Lemma 1. This estimate is optimal for solutions with asymptotics of the form (9)-(10) if we assume that $`0<v<1`$ and that $`\varphi `$ is not identically zero.

###### Theorem 1.

Consider a solution to (2)-(3). Assume that $`\tau _02`$,

(12)
$$1P(\tau _0,\theta )\tau _01,\gamma \frac{1}{\tau _0}P(\tau _0,\theta )1\gamma $$

for $`\theta S^1`$ and some $`\gamma >0`$, and that $`F(\tau _0)(\gamma \alpha )^2`$ for some $`\alpha >0`$, $`\alpha <\gamma `$. Then there are $`v,\varphi ,q,rC^{\mathrm{}}(S^1,)`$, and for all non-negative integers $`k`$, polynomials $`\mathrm{\Xi }_{i,k}`$, $`i=1,\mathrm{},4`$ in $`\tau `$, where

(13)
$$0<\alpha v1\alpha <1$$

on $`S^1`$ such that

(14)
$$P_\tau v_{C^k(S^1,)}\mathrm{\Xi }_{1,k}\mathrm{exp}[\alpha \tau ],$$

(15)
$$Pp_{C^k(S^1,)}\mathrm{\Xi }_{2,k}\mathrm{exp}[\alpha \tau ],$$

(16)
$$e^{2p}Q_\tau r_{C^k(S^1,)}\mathrm{\Xi }_{3,k}\mathrm{exp}[\alpha \tau ],$$

and

(17)
$$e^{2p}(Qq)+\frac{r}{2v}_{C^k(S^1,)}\mathrm{\Xi }_{4,k}\mathrm{exp}[\alpha \tau ]$$

for all $`\tau [\tau _0,\mathrm{})`$, where $`p=v\tau +\varphi `$.

Remark. The result obtained is the same as in , but the conditions are weaker. We wish to emphasize that the proof does not rely on  in any essential way. We only need two pages of estimates from that paper in order to complete the argument. Thus, the rather intricate arguments presented in  are avoided.

The proof is to be found at the end of Section 2. Under the conditions of the theorem, we thus obtain the expansions (9)-(10). Note the condition (12) which requires that $`P1`$. This is a rather unnatural condition which one would like to get rid of. It turns out that the reason one has to impose this condition is due to a bad choice of representative of hyperbolic space. Consider the asymptotic behaviour (9)-(10). Under the map (8), one gets the following picture in the upper half plane. The $`x`$-coordinate converges for every fixed $`\theta `$ and the $`y`$-coordinate tends to zero. In other words, for a fixed $`\theta `$, the solution goes to the boundary along a geodesic, at least approximately. Let us consider the same picture in the disc model. Let $`D`$ be the open unit disc in the complex plane and let

(18)
$$g_D=\frac{4(dx^2+dy^2)}{(1x^2y^2)^2}$$

be a metric on it. In complex notation, we have the isometry

(19)
$$\varphi _{HD}(z)=\frac{zi}{z+i}$$

from the upper half plane to the disc model. Note that the inverse is given by $`\varphi _{HD}^1(w)=i(1+w)/(1w)`$. Say now that we have a solution tending to the boundary in the disc model. For some spatial point, the solution tends to $`1`$, and for neighbouring points, the solution goes to the boundary a little bit below and a little bit above in the complex plane. If one translates this to the $`PQ`$-variables, one gets a very violent behaviour which is completely unrelated to the geometry of the wave map problem. For this reason, we consider the equations in the disc model. We have

(20)
$$_\tau \left(\frac{z_\tau }{(1|z|^2)^2}\right)e^{2\tau }_\theta \left(\frac{z_\theta }{(1|z|^2)^2}\right)=\frac{2z}{(1|z|^2)^3}[|z_\tau |^2e^{2\tau }|z_\theta |^2].$$

We will use the notation

(21)
$$|z_\tau |_D^2=\frac{4}{(1|z|^2)^2}|z_\tau |^2$$

and similarly for $`z_\theta `$, $`z_\tau +e^\tau z_\theta `$ etc. Let

(22)
$$\rho =\mathrm{ln}\frac{1+|z|}{1|z|}.$$

This is the distance from the origin to the point $`z`$ with respect to the hyperbolic metric. Note that $`\rho ^2`$ and $`\rho /|z|`$ are smooth functions on the open unit disc.

###### Theorem 2.

Consider a solution $`z`$ to (20). Assume that there is an $`0<\alpha <1`$ such that at $`\tau _0`$,

(23)
$$\frac{1}{2}\underset{\theta }{sup}|z_\tau +e^\tau z_\theta |_D^2+\frac{1}{2}\underset{\theta }{sup}|z_\tau e^\tau z_\theta |_D^2(1\alpha )^2.$$

Then there is a $`vC^0(S^1,^2)`$ with $`|v|1\alpha `$ and a $`C`$ and a $`T`$ such that for $`\tau T`$,

$$\frac{1}{\tau }\frac{z(\tau ,)}{|z(\tau ,)|}\rho (\tau ,)v_{C^0(S^1,^2)}+\frac{2z_\tau (\tau ,)}{1|z(\tau ,)|^2}v_{C^0(S^1,^2)}$$

$$+e^\tau \frac{2z_\theta (\tau ,)}{1|z(\tau ,)|^2}_{C^0(S^1,^2)}C\tau ^1.$$

Remark. Note that if $`v(\theta _0)=0`$, then the above estimate implies that $`z(\tau ,\theta _0)`$ remains bounded as $`\tau \mathrm{}`$. In the proof, we show that the left hand side of (23) decreases as $`\tau `$ increases.

The proof is to be found at the end of Section 3. As noted in Lemma 5, one gets the same conclusions if one replaces the condition (23) with the condition that there is a $`T`$ such that $`\rho (\tau ,\theta )\tau 2`$ for $`\tau T`$. Of course, one would like to have smooth expansions. Note that this can be done in all generality in the polarized case, see . We are not able to prove that there are smooth expansions if $`|v|=0`$, but all other cases can be dealt with in the following sense.

###### Theorem 3.

Consider a solution to (20). Let $`\theta _0S^1`$ be a fixed angle, and assume that there is a $`T`$ and a $`0<\gamma <1`$ such that

$$\frac{\rho (\tau ,\theta _0)}{\tau }1\gamma $$

for all $`\tau T`$. Assume furthermore that there is a sequence $`\tau _k\mathrm{}`$ such that

$$\frac{\rho (\tau _k,\theta _0)}{\tau _k}\gamma $$

Then there is an $`\eta >0`$ and an isometry $`\varphi `$ from $`(D,g_D)`$ to $`(^2,g_R)`$ such that if $`(Q,P)=\varphi z`$, then the conclusions of Theorem 1 hold if we replace $`S^1`$ with $`I_\eta `$, where $`I_\eta =(\theta _0\eta ,\theta _0+\eta )`$.

The theorem follows from Lemma 6 and Lemma 7. If we skip the condition that $`\rho (\tau _k,\theta _0)/\tau _k\gamma `$, we still get conclusions similar to those of Theorem 2, cf. Lemma 7.

When working on this article, the author was made aware of the fact that similar results had been obtained by Chae and Chruściel. Since their methods are quite different from ours, it is however our hope that the reader will find both articles interesting.

## 2. Estimates

The purpose of this section is to prove Theorem 1. We begin by describing the decay properties of the function $`F`$ defined in (11).

###### Lemma 1.

Consider a solution to (2)-(3). Assume that $`\tau _2\tau _12`$ and that

$$1P(\tau ,\theta )\tau 1$$

for $`(\tau ,\theta )[\tau _1,\tau _2]\times S^1`$. Then if $`F`$ is defined in (11),

$$F(\tau )F(\tau _1)\left(\frac{\tau _1}{\tau }\right)^2$$

for all $`\tau [\tau _1,\tau _2]`$.

Proof. Let

$$𝒜_1=\frac{1}{2}e^\tau (P_\tau \frac{1}{\tau }P+e^\tau P_\theta )^2,𝒜_2=\frac{1}{2}e^\tau e^{2P}(Q_\tau +e^\tau Q_\theta )^2,$$

$$_1=\frac{1}{2}e^\tau (P_\tau \frac{1}{\tau }Pe^\tau P_\theta )^2,_2=\frac{1}{2}e^\tau e^{2P}(Q_\tau e^\tau Q_\theta )^2$$

and

(24)
$$𝒜=𝒜_1+𝒜_2,=_1+_2.$$

One can compute that

$$(_\tau e^\tau _\theta )𝒜_1=(\frac{1}{2}\frac{1}{\tau })e^\tau [(P_\tau \frac{1}{\tau }P)^2+e^{2\tau }P_\theta ^2](1\frac{2}{\tau })e^\tau P_\theta ^2$$

$$+e^{2P+\tau }(Q_\tau ^2e^{2\tau }Q_\theta ^2)(P_\tau \frac{1}{\tau }P+e^\tau P_\theta )$$

and that

(25)
$$(_\tau e^\tau _\theta )𝒜_2=\frac{1}{2}e^{2P+\tau }[Q_\tau ^2+e^{2\tau }Q_\theta ^2]$$

$$e^{2P+\tau }(P_\tau +e^\tau P_\theta )(Q_\tau ^2e^{2\tau }Q_\theta ^2)e^{2P\tau }Q_\theta ^2.$$

Thus

$$(_\tau e^\tau _\theta )𝒜=(\frac{1}{2}\frac{1}{\tau })e^\tau [(P_\tau \frac{1}{\tau }P)^2+e^{2\tau }P_\theta ^2](1\frac{2}{\tau })e^\tau P_\theta ^2$$

$$\frac{P}{\tau }e^{2P+\tau }(Q_\tau ^2e^{2\tau }Q_\theta ^2)+\frac{1}{2}e^{2P+\tau }[Q_\tau ^2+e^{2\tau }Q_\theta ^2]e^{2P\tau }Q_\theta ^2.$$

Assuming $`1P\tau 1`$ and $`\tau 2`$, we thus get

$$(_\tau e^\tau _\theta )𝒜(\frac{1}{2}\frac{1}{\tau })(𝒜+).$$

Similarly, if $`1P\tau 1`$ and $`\tau 2`$,

$$(_\tau +e^\tau _\theta )(\frac{1}{2}\frac{1}{\tau })(𝒜+).$$

Let us write down the analogue of (25) for future reference.

(26)
$$(_\tau +e^\tau _\theta )_2=\frac{1}{2}e^{2P+\tau }[Q_\tau ^2+e^{2\tau }Q_\theta ^2]$$

$$e^{2P+\tau }(P_\tau e^\tau P_\theta )(Q_\tau ^2e^{2\tau }Q_\theta ^2)e^{2P\tau }Q_\theta ^2.$$

Let

$$F_1(u,\theta )=𝒜(u,\theta +e^u),F_2(u,\theta )=(u,\theta e^u),$$

(27)
$$\widehat{F}_i(u)=\underset{\theta S^1}{sup}F_i(u,\theta ),\widehat{F}(u)=\widehat{F}_1(u)+\widehat{F}_2(u).$$

We get

$$F_1(u_2,\theta )=F_1(u_1,\theta )+_{u_1}^{u_2}[(_ue^u_\theta )𝒜](u,\theta +e^u)𝑑u$$

$$\widehat{F}_1(u_1)+_{u_1}^{u_2}(\frac{1}{2}\frac{1}{u})\widehat{F}(u)𝑑u,$$

assuming $`1P(\tau ,\theta )\tau 1`$ and $`\tau 2`$ for $`(\tau ,\theta )[u_1,u_2]\times S^1`$. Take the sup norm of this and add it to a similar estimate for $`\widehat{F}_2`$ in order to obtain

$$\widehat{F}(u_2)\widehat{F}(u_1)+_{u_1}^{u_2}(1\frac{2}{u})\widehat{F}(u)𝑑u.$$

Grönwall’s lemma implies

$$e^{u_2}\widehat{F}(u_2)e^{u_1}\widehat{F}(u_1)\left(\frac{u_1}{u_2}\right)^2.$$

The lemma follows. $`\mathrm{}`$

###### Corollary 1.

Consider a solution to (2)-(3). Assume that $`\tau _02`$,

$$1P(\tau _0,\theta )\tau _01,\gamma \frac{1}{\tau _0}P(\tau _0,\theta )1\gamma $$

for $`\theta S^1`$ and some $`\gamma >0`$, and that $`F(\tau _0)(\gamma \alpha )^2`$ for some $`\alpha >0`$, $`\alpha <\gamma `$. Then

$$F(\tau )F(\tau _0)\left(\frac{\tau _0}{\tau }\right)^2$$

for all $`\tau \tau _0`$. Furthermore, there is a function $`vC^0(S^1,)`$ and a constant $`0<C<\mathrm{}`$ such that

$$P_\tau (\tau ,)v_{C^0(S^1,)}\frac{C}{\tau }$$

for $`\tau \tau _0`$ and $`\alpha v(\theta )1\alpha `$ for all $`\theta S^1`$.

Proof. Consider the set

$$𝒮=\{\tau \tau _0:s[\tau _0,\tau ]1P(s,\theta )s1\theta S^1\}.$$

This set is closed and connected by definition. From the conditions of the lemma, it is clear that it is non-empty. We wish to prove that it is open (in the subspace topology on $`[\tau _0,\mathrm{})`$). Assume that $`\tau _1𝒮`$. Let us estimate, using Lemma 1,

(28)
$$|\frac{P(\tau _1,\theta )}{\tau _1}\frac{P(\tau _0,\theta )}{\tau _0}|=|_{\tau _0}^{\tau _1}\frac{1}{\tau }[P_\tau (\tau ,\theta )\frac{1}{\tau }P(\tau ,\theta )]𝑑\tau |_{\tau _0}^{\tau _1}\frac{F^{1/2}(\tau )}{\tau }𝑑\tau $$

$$(\gamma \alpha )_{\tau _0}^{\tau _1}\frac{\tau _0}{\tau ^2}𝑑\tau (\gamma \alpha )(1\frac{\tau _0}{\tau _1}).$$

Using Lemma 1 once again we conclude that

$$|P_\tau (\tau _1,\theta )\frac{P(\tau _1,\theta )}{\tau _1}|(\gamma \alpha )\frac{\tau _0}{\tau _1}.$$

Combining these two estimates, we get

$$|P_\tau (\tau _1,\theta )\frac{P(\tau _0,\theta )}{\tau _0}|\gamma \alpha .$$

Since $`\gamma P(\tau _0,\theta )/\tau _01\gamma `$ for all $`\theta S^1`$, $`\alpha P_\tau (\tau _1,\theta )1\alpha `$ for all $`\theta S^1`$. This implies that there is a set of the form $`[\tau _1,\tau _1+ϵ)`$ with $`ϵ>0`$ such that $`1P(s,\theta )s1`$ for all $`s`$ in this set and $`\theta S^1`$. Consequently, $`𝒮`$ is open, so that $`𝒮=[\tau _0,\mathrm{})`$. An estimate similar to that of (28) yields the conclusion that for $`\tau _2\tau _1\tau _0`$

$$\frac{P(\tau _2,)}{\tau _2}\frac{P(\tau _1,)}{\tau _1}_{C^0(S^1,)}\frac{C}{\tau _1}.$$

Thus there is a $`vC^0(S^1,)`$ such that

$$\frac{P(\tau ,)}{\tau }v_{C^0(S^1,)}\frac{C}{\tau }.$$

Letting $`\tau _1`$ tend to infinity in (28), we get the conclusion that $`\alpha v(\theta )1\alpha `$ for all $`\theta S^1`$. Combining this with the fact that $`F`$ decays like $`\tau ^2`$, we get the conclusions of the corollary. $`\mathrm{}`$

###### Lemma 2.

Consider a solution to (2)-(3) with initial data satisfying the conditions in the statement of Corollary 1. Then

$$e^{2P}[Q_\tau ^2(\tau ,)+e^{2\tau }Q_\theta ^2(\tau ,)]_{C^0(S^1,)}Ce^{\alpha \tau }.$$

Proof. Consider (25) and (26). By Corollary 1 we conclude that $`e^\tau P_\theta `$ converges to zero in the sup norm. Since we also know that $`P_\tau `$ converges to $`v`$, we conclude that for $`\tau `$ great enough,

$$\frac{\alpha }{2}P_\tau \pm e^\tau P_\theta 1\frac{\alpha }{2}.$$

Combining this with (25) and (26) we obtain

$$(_\tau e^\tau _\theta )𝒜_2(\frac{1}{2}\frac{\alpha }{2})(𝒜_2+_2),(_\tau +e^\tau _\theta )_2(\frac{1}{2}\frac{\alpha }{2})(𝒜_2+_2).$$

The remaining part of the argument is similar to the end of the proof of Lemma 1. $`\mathrm{}`$

Consider the energies

$$_k=\frac{1}{2}_{S^1}[(_\theta ^k_\tau P)^2+e^{2\tau }(_\theta ^{k+1}P)^2]𝑑\theta $$

and

$$E_k=\frac{1}{2}_{S^1}e^{2P}[(_\theta ^k_\tau Q)^2+e^{2\tau }(_\theta ^{k+1}Q)^2]𝑑\theta .$$

We will also be interested in the energies

(29)
$$H_k=1+_k+E_k+\frac{1}{2}e^{\alpha \tau }_{S^1}(_\theta ^kP)^2𝑑\theta .$$

Note that

(30)
$$\frac{d_k}{d\tau }_{S^1}_\theta ^k[e^{2P}(Q_\tau ^2e^{2\tau }Q_\theta ^2)]_\theta ^k_\tau Pd\theta .$$

Furthermore, since we can assume $`\alpha /2P_\tau 1\alpha /2`$ for $`\tau `$ large enough, one can estimate

(31)
$$\frac{dE_k}{d\tau }\alpha E_k+_{S^1}f_k_\theta ^k_\tau Qd\theta $$

for $`\tau `$ large enough, where

(32)
$$f_k=_\tau (e^{2P}_\theta ^k_\tau Q)e^{2\tau }_\theta (e^{2P}_\theta ^{k+1}Q).$$

Note that $`f_0=0`$ and that

(33)
$$f_{k+1}=_\theta f_k2P_\theta f_k2P_{\theta \tau }e^{2P}_\theta ^k_\tau Q+2P_{\theta \theta }e^{2P2\tau }_\theta ^{k+1}Q.$$

The details of the derivations of the last three equations can be found in the proof of Lemma 5.2 of .

###### Lemma 3.

Consider a solution to (2)-(3) with initial data satisfying the conditions in the statement of Corollary 1. Then $`H_1`$ defined by (29) is bounded to the future.

Proof. We wish to prove this by proving an estimate of the form

$$\frac{dH_1}{d\tau }Ce^{\beta \tau }H_1,$$

where $`\beta >0`$. Let us prove this with $`H_1`$ on the left hand side replaced by the constituents of $`H_1`$. By (30) we have

$$\frac{d_1}{d\tau }_{S^1}_\theta [e^{2P}(Q_\tau ^2e^{2\tau }Q_\theta ^2)]P_{\tau \theta }d\theta \sqrt{2}_\theta [e^{2P}(Q_\tau ^2e^{2\tau }Q_\theta ^2)]_{L^2(S^1,)}_1^{1/2}.$$

Let us estimate the $`L^2`$-norm appearing on the right hand side. There are two cases to deal with. Estimate, using Lemma 2,

$$P_\theta e^{2P}(Q_\tau ^2e^{2\tau }Q_\theta ^2)_{L^2(S^1,)}Ce^{\alpha \tau }e^{\alpha \tau /2}H_1^{1/2},$$

which is of the desired form. The other case can be estimated

$$e^{2P}(Q_\tau Q_{\tau \theta }e^{2\tau }Q_\theta Q_{\theta \theta })_{L^2(S^1,)}Ce^{\alpha \tau /2}E_1^{1/2},$$

using Lemma 2. By (31)-(33) and the fact that $`f_0=0`$, we get

$$\frac{dE_1}{d\tau }_{S^1}[2P_{\tau \theta }e^{2P}Q_\tau Q_{\tau \theta }+2e^{2P2\tau }P_{\theta \theta }Q_\theta Q_{\tau \theta }]𝑑\theta .$$

Thus we can use Lemma 2 in order to obtain

$$\frac{dE_1}{d\tau }Ce^{\alpha \tau /2}(E_1+_1).$$

Finally

$$\frac{d}{d\tau }[\frac{1}{2}e^{\alpha \tau }_{S^1}P_\theta ^2𝑑\theta ]=\frac{\alpha }{2}e^{\alpha \tau }_{S^1}P_\theta ^2𝑑\theta +e^{\alpha \tau }_{S^1}P_\theta P_{\tau \theta }𝑑\theta 2e^{\alpha \tau /2}H_1^{1/2}_1^{1/2}.$$

The lemma follows. $`\mathrm{}`$

###### Corollary 2.

Consider a solution to (2)-(3) with initial data satisfying the conditions in the statement of Corollary 1. Then

(34)
$$E_1(\tau )C(1+\tau ^2)e^{\alpha \tau }$$

for all $`\tau \tau _0`$.

Proof. By (31)-(33) and the fact that $`f_0=0`$, we get

$$\frac{dE_1}{d\tau }\alpha E_1+C\mathrm{exp}[\alpha \tau /2]E_1^{1/2},$$

where we have also used Lemma 2 and the fact that $`H_1`$ is bounded. This yields the conclusion of the corollary. $`\mathrm{}`$

###### Lemma 4.

Consider a solution to (2)-(3) with initial data satisfying the conditions in the statement of Corollary 1. Then for every $`k`$, there is a constant $`C_k`$ and a polynomial $`𝒫_k`$ in $`\tau `$ such that for $`\tau 0`$

$$H_kC_k,E_k𝒫_k(\tau )e^{\alpha \tau }.$$

Proof. Let us make the inductive assumption that

$$H_iC_i,E_i𝒫_i(\tau )e^{\alpha \tau }$$

for $`i=1,\mathrm{},k`$, where the $`C_i`$ are constants and the $`𝒫_i`$ are polynomials. By Lemma 3 and (34), we know this to be true for $`k=1`$. Note that if we let

$$_l=\frac{1}{2}_{S^1}(_\theta ^{l+1}P)^2𝑑\theta ,$$

we have

$$\frac{d_l}{d\tau }=_{S^1}_\theta ^{l+1}P_\theta ^{l+1}_\tau Pd\theta 2_l^{1/2}_{l+1}^{1/2},$$

whence

$$_lC_l(1+\tau ^2)$$

for $`l=0,\mathrm{},k1`$. Consequently, we have, using the definition of $`H_k`$,

(35)
$$_\theta ^lP_{C^0(S^1,)}C_l(1+\tau ^2)^{1/2},_\theta ^kP_{C^0(S^1,)}C_ke^{\alpha \tau /2}H_{k+1}^{1/2},$$

(36)
$$_\theta ^kP_{L^2(S^1,)}C_k(1+\tau ^2)^{1/2},_\theta ^{k+1}P_{L^2(S^1,)}C_ke^{\alpha \tau /2}H_{k+1}^{1/2},$$

where $`l=1,\mathrm{},k1`$. Note that if $`l1`$, $`_\theta ^l_\tau Q`$ has to have a zero on the circle since its average over the circle is zero. Thus

(37)
$$e^P_\theta ^l_\tau Q_{C^0(S^1,)}_{S^1}|P_\theta e^P_\theta ^l_\tau Q+e^P_\theta ^{l+1}_\tau Q|𝑑\theta $$

$$𝒬_le^{\alpha \tau /2}+2\pi ^{1/2}E_{l+1}^{1/2},$$

for $`l=1,\mathrm{},k`$, where $`𝒬_l`$ is a polynomial in $`\tau `$ and we have used the induction hypothesis and the first inequality in (36) with $`k=1`$, which holds due to Lemma 3. If $`lk1`$, we thus get exponential decay of the left hand side, due to the induction hypothesis. A similar estimate holds if we replace $`Q_\tau `$ with $`e^\tau Q_\theta `$.

The idea of proof is the same as in Lemma 3. Consider

$$\frac{d_{k+1}}{d\tau }_{S^1}_\theta ^{k+1}[e^{2P}(Q_\tau ^2e^{2\tau }Q_\theta ^2)]_\theta ^{k+1}_\tau Pd\theta .$$

By Hölder’s inequality, we need to estimate expressions of the form

(38)
$$_\theta ^{l_1}P\mathrm{}_\theta ^{l_o}Pe^{2P}_\theta ^m_\tau Q_\theta ^n_\tau Q_{L^2(S^1,)}$$

where $`l_1+\mathrm{}+l_o+m+n=k+1`$, and the analogous expressions with $`Q_\tau `$ replaced by $`e^\tau Q_\theta `$. We wish to estimate this by $`H_{k+1}^{1/2}`$ times something that decays exponentially. If the biggest $`l_i`$ is $`k+1`$, we can use the second inequality in (36) and Lemma 2. If there are two $`l_i`$ which equal $`k`$, then $`k=1`$ and we can use Lemma 2, the second inequality in (35) and the first inequality in (36) to obtain the desired estimate. If there is only one $`l_i`$ equal to $`k`$ and the rest are smaller, we can take out all factors $`_\theta ^{l_j}P`$ using the first or the second of (35). The remaining terms are then estimated using Lemma 2 or (34). If all the $`l_i`$ are less than or equal to $`k1`$, we can take out the derivatives of $`P`$ in the sup norm with a polynomial bound in $`\tau `$, cf. the first inequality in (35). Thus, we need not concern ourselves with derivatives hitting $`P`$ in what follows. If one of $`m`$ and $`n`$ are less than or equal to $`k1`$, we can take out the corresponding factor in the sup norm, using (37) and the fact that $`E_{l+1}`$ decays exponentially for $`lk1`$. The problem which remains is the case where $`m=n=k`$. Then $`k=1`$ and we can use (34) and (37) to take out one factor in the sup norm and estimate the remaining $`L^2`$ norm by $`E_1^{1/2}`$. The argument for expressions of the form (38) where one has replaced $`Q_\tau `$ with $`e^\tau Q_\theta `$ is similar.

Consider (31)-(33). Since $`f_0=0`$, we inductively get the conclusion that $`f_{k+1}`$ is a sum of terms of the form

(39)
$$_\theta ^{m+1}_\tau Pe^{2P}_\theta ^l_\tau Q,$$

where $`m+l=k`$ and terms where one replaces $`P_\tau `$ with $`e^\tau P_\theta `$ and $`Q_\tau `$ with $`e^\tau Q_\theta `$. Note that the effect of the operator $`_\theta 2P_\theta `$ on an expression of the form $`e^{2P}g`$ is to only differentiate $`g`$. This is the reason for (39). Discarding the first term in (31), we see that the estimate we wish to achieve is

$$e^Pf_{k+1}_{L^2(S^1,)}Ce^{\beta \tau }H_{k+1}^{1/2}.$$

Due to the form (39), the relevant things to estimate are thus

(40)
$$_\theta ^{m+1}_\tau Pe^P_\theta ^l_\tau Q_{L^2(S^1,)},$$

where $`l+m=k`$. If $`lk1`$, we use (37) to estimate $`e^P_\theta ^l_\tau Q`$ in the sup norm. If $`l=k`$, we take out $`P_{\tau \theta }`$ in the sup norm and use the induction hypothesis to conclude that what remains decays exponentially. If $`k=1`$, we can estimate $`P_{\tau \theta }`$ by $`H_{k+1}^{1/2}`$. Otherwise, it is bounded. The arguments for the terms where one replaces $`P_\tau `$ with $`e^\tau P_\theta `$ and $`Q_\tau `$ with $`e^\tau Q_\theta `$ are similar.

Finally

$$\frac{d}{d\tau }[\frac{1}{2}e^{\alpha \tau }_{S^1}(_\theta ^{k+1}P)^2𝑑\theta ]e^{\alpha \tau }_{S^1}_\theta ^{k+1}P_\theta ^{k+1}_\tau Pd\theta 2e^{\alpha \tau /2}H_{k+1}.$$

In consequence, we have the desired inequality for $`H_{k+1}^{}`$ and we conclude that $`H_{k+1}`$ is bounded. In order to complete the induction, we need to prove the decay of $`E_{k+1}`$. Due to (31), we only need to prove that expressions of the form (40) can be estimated by a polynomial times $`e^{\alpha \tau /2}`$ (and the analogous expressions where $`_\tau `$ is replaced with $`e^\tau _\theta `$). If $`lk1`$, we can use (37) and the boundedness of $`H_{k+1}`$. Since $`P_{\tau \theta }`$ is bounded in the sup norm due to the boundedness of $`H_2`$, and since $`E_k`$ has the desired decay, we can also deal with the case $`l=k`$. The lemma follows. $`\mathrm{}`$

Proof of Theorem 1. The conditions of Corollary 1 are satisfied. By Lemma 4, we conclude that $`_k`$ is bounded and that $`E_k`$ decays as $`e^{\alpha \tau }`$ times a polynomial for all $`k`$. Furthermore, for all $`k`$, there is a polynomial $`𝒬_k`$ in $`\tau `$ such that for all $`k`$

$$e^P_\theta ^k_\tau Q_{C^0(S^1,)}+e^{P\tau }_\theta ^{k+1}Q_{C^0(S^1,)}𝒬_k(\tau )e^{\alpha \tau /2}.$$

This follows from (37) and Lemma 4. When one has these estimates, it is not so difficult to combine them with the equations in order to get the conclusions of the theorem. However, the arguments require almost two pages, so we refer the reader to pp. 23-25 of  for the details. $`\mathrm{}`$

## 3. The disc model

We begin by proving a result analogous to Lemma 1 in the disc model.

###### Lemma 5.

Let $`z`$ be a solution to (20) and assume that there is a $`T`$ such that $`\rho (\tau ,\theta )\tau 2`$ for $`(\tau ,\theta )[T,\mathrm{})\times S^1`$, where $`\rho `$ is defined in (22). Then there is a $`vC^0(S^1,^2)`$ with $`|v|1`$ and a $`C`$ such that for $`\tau T`$,

$$\frac{1}{\tau }\frac{z(\tau ,)}{|z(\tau ,)|}\rho (\tau ,)v_{C^0(S^1,^2)}+\frac{2z_\tau (\tau ,)}{1|z(\tau ,)|^2}v_{C^0(S^1,^2)}$$

$$+e^\tau \frac{2z_\theta (\tau ,)}{1|z(\tau ,)|^2}_{C^0(S^1,^2)}C\tau ^1.$$

Proof. If $`z`$ is a solution of (20), let

(41)
$$𝒜_1=\frac{1}{2}e^\tau |z_\tau +e^\tau z_\theta |_D^2,𝒜_2=\frac{1}{2}e^\tau |z_\tau e^\tau z_\theta |_D^2,$$

with notation as in (21), and

$$𝒞_1=\frac{1}{2\tau }e^\tau (_\tau +e^\tau _\theta )(\rho )^2+\frac{1}{2}e^\tau \frac{\rho ^2}{\tau ^2},𝒞_2=\frac{1}{2\tau }e^\tau (_\tau e^\tau _\theta )(\rho )^2+\frac{1}{2}e^\tau \frac{\rho ^2}{\tau ^2}.$$

Note that $`𝒞_1`$ and $`𝒞_2`$ are smooth functions. If we let

$$f=\frac{1}{2\tau }(1|z|^2)\frac{z}{|z|}\rho ,$$

which is a smooth function, then

(42)
$$_1=𝒜_1+𝒞_1=\frac{1}{2}e^\tau |z_\tau f+e^\tau z_\theta |_D^2$$

and similarly for $`_2=𝒜_2+𝒞_2`$. Consequently $`_1`$ and $`_2`$ are non-negative. We have

(43)
$$(_\tau e^\tau _\theta )𝒜_1=(_\tau +e^\tau _\theta )𝒜_2=\frac{1}{2}e^\tau [|z_\tau |_D^2e^{2\tau }|z_\theta |_D^2].$$

Let us compute

$$(_\tau e^\tau _\theta )𝒞_1=\frac{1}{\tau ^2}e^\tau _\tau (\rho ^2)\frac{1}{2\tau }e^\tau _\tau (\rho ^2)+\frac{1}{2}e^\tau \frac{\rho ^2}{\tau ^2}\frac{1}{\tau }e^\tau \frac{\rho ^2}{\tau ^2}$$

$$\frac{1}{2\tau }e^\tau (_\tau ^2e^{2\tau }_\theta ^2)(\rho ^2).$$

If $`|z|>0`$, we can compute

(44)
$$\frac{1}{2}(_\tau ^2e^{2\tau }_\theta ^2)(\rho ^2)=(\rho _\tau ^2e^{2\tau }\rho _\theta ^2)+\rho (\rho _{\tau \tau }e^{2\tau }\rho _{\theta \theta }),$$

and that

$$\mathrm{sinh}\rho (\rho _{\tau \tau }e^{2\tau }\rho _{\theta \theta })=\mathrm{cosh}\rho [|z_\tau |_D^2e^{2\tau }|z_\theta |_D^2\rho _\tau ^2+e^{2\tau }\rho _\theta ^2].$$

Observe that

(45)
$$\rho _\tau ^2+\mathrm{sinh}^2\rho |_\tau (\frac{z}{|z|})|^2=|z_\tau |_D^2$$

and similarly for the $`\theta `$ derivatives. Consequently

$$|z_\tau |_D^2\rho _\tau ^20.$$

Since $`\rho \mathrm{cosh}\rho \mathrm{sinh}\rho `$, we conclude that

$$\rho (\rho _{\tau \tau }e^{2\tau }\rho _{\theta \theta })|z_\tau |_D^2+\rho _\tau ^2e^{2\tau }\rho _\theta ^2+\frac{\rho \mathrm{cosh}\rho }{\mathrm{sinh}\rho }e^{2\tau }|z_\theta |_D^2$$

Combining this with (44), we get

$$\frac{1}{2}(_\tau ^2e^{2\tau }_\theta ^2)(\rho ^2)|z_\tau |_D^2+\frac{\rho \mathrm{cosh}\rho }{\mathrm{sinh}\rho }e^{2\tau }|z_\theta |_D^2.$$

Since $`\rho \mathrm{cosh}\rho /\mathrm{sinh}\rho \tau 1`$ if $`\rho \tau 2`$, we get

$$\frac{1}{2\tau }e^\tau (_\tau ^2e^{2\tau }_\theta ^2)(\rho ^2)e^\tau |z_\theta |_D^2\frac{1}{\tau }e^\tau [|z_\tau |_D^2+e^{2\tau }|z_\theta |_D^2].$$

We conclude that if $`\rho \tau 2`$, then

(46)
$$(_\tau e^\tau _\theta )_1(\frac{1}{2}\frac{1}{\tau })e^\tau [|z_\tau |_D^2+e^{2\tau }|z_\theta |_D^2]+\frac{1}{\tau ^2}e^\tau _\tau (\rho ^2)$$

$$\frac{1}{2\tau }e^\tau _\tau (\rho ^2)+\frac{1}{2}e^\tau \frac{\rho ^2}{\tau ^2}\frac{1}{\tau }e^\tau \frac{\rho ^2}{\tau ^2}=(\frac{1}{2}\frac{1}{\tau })(_1+_2).$$

Similarly,

(47)
$$(_\tau +e^\tau _\theta )_2(\frac{1}{2}\frac{1}{\tau })(_1+_2).$$

The above derivation was carried out under the assumption that $`|z|>0`$. However, if $`z(\tau ,\theta )=0`$, then $`\rho ^2`$ can be replaced with $`4|z|^2`$ when calculating second derivatives. Thus, at the point $`(\tau ,\theta )`$, we have

$$(_\tau ^2e^{2\tau }_\theta ^2)(\rho ^2)=8(|z_\tau |^2e^{2\tau }|z_\theta |^2).$$

Using this observation, one can see that (46) and (47) hold regardless of whether $`|z|=0`$ or not. Let

$$G(\tau )=e^\tau \underset{\theta }{sup}_1(\tau ,\theta )+e^\tau \underset{\theta }{sup}_2(\tau ,\theta ),$$

and assume that there is a $`T`$ such that $`\rho (\tau ,\theta )\tau 2`$ for all $`(\tau ,\theta )[T,\mathrm{})\times S^1`$. Then an argument similar to the proof of Lemma 1 shows that

$$G(\tau )G(\tau _0)\left(\frac{\tau _0}{\tau }\right)^2$$

if $`\tau \tau _0T`$. Assuming $`|z|>0`$, we can use (45) in order to obtain

$$e^\tau _1+e^\tau _2=(\rho _\tau \frac{1}{\tau }\rho )^2+\mathrm{sinh}^2\rho |_\tau (\frac{z}{|z|})|^2+e^{2\tau }|z_\theta |_D^2.$$

Consequently, there is a constant $`C`$ such that for $`\tau T`$

$$(\rho _\tau \frac{1}{\tau }\rho )^2+\mathrm{sinh}^2\rho |_\tau (\frac{z}{|z|})|^2C\tau ^2$$

assuming $`|z|>0`$. Define $`g`$ by

$$g=\frac{1}{\tau }\frac{z}{|z|}\rho .$$

Let us compute, assuming $`|z|>0`$,

$$_\tau g=\frac{z}{|z|}\frac{1}{\tau }(\rho _\tau \frac{1}{\tau }\rho )+\frac{1}{\tau }_\tau (\frac{z}{|z|})\rho .$$

Since $`\rho /\mathrm{sinh}\rho `$ is bounded, we get the conclusion that

(48)
$$|_\tau g|C\tau ^2$$

for $`\tau T`$, assuming $`|z|>0`$. If $`z=0`$, we get

$$_\tau g=\frac{2}{\tau }z_\tau .$$

However, in this case

$$e^\tau _1+e^\tau _2=4|z_\tau |^2+4e^{2\tau }|z_\theta |^2.$$

Consequently (48) holds regardless of whether $`z`$ is zero or not. If $`\tau _2\tau _1T`$, we thus get the conclusion that

$$g(\tau _2,)g(\tau _1,)_{C^0(S^1,)}C\tau _1^1.$$

In other words, there is a $`vC^0(S^1,^2)`$ such that

$$g(\tau ,)v_{C^0(S^1,^2)}C\tau ^1.$$

Note that $`|v|1`$, since $`|g|1`$ for $`\tau T`$. Since

$$e^\tau _1+e^\tau _2=|\frac{2z_\tau }{1|z|^2}g|^2+e^{2\tau }|z_\theta |_D^2,$$

the lemma follows. $`\mathrm{}`$

Proof of Theorem 2. Consider $`𝒜_1`$ and $`𝒜_2`$ defined in (41). Due to (43), we have

$$(_\tau e^\tau _\theta )𝒜_1=(_\tau +e^\tau _\theta )𝒜_2\frac{1}{2}(𝒜_1+𝒜_2).$$

This can be used to conclude that the object appearing on the left hand side of (23) is monotonically decaying to the future. The argument is similar to the proof of Lemma 1. Consequently, $`|z_\tau |_D1\alpha `$ to the future. Due to (45), this means that $`\rho (\tau ,\theta )(1\alpha )\tau +C`$ for $`\tau \tau _0`$. At late enough times, the conditions of Lemma 5 are thus fulfilled. The theorem follows. $`\mathrm{}`$

## 4. Local results

Here we wish to state some results that are local in the spatial coordinate. We also wish to relate the pictures in the different models of hyperbolic space. Note that we have the three models of hyperbolic space (6), (7) and (18) which are related by the isometries (8) and (19).

###### Lemma 6.

Consider a solution to (20). Let $`\theta _0S^1`$ and assume that there are $`\gamma ,ϵ>0`$, $`T`$ and $`vC^0(I_ϵ,^2)`$, where $`I_ϵ=(\theta _0ϵ,\theta _0+ϵ)`$, $`2\gamma |v|12\gamma `$ in $`I_ϵ`$ and

(49)
$$\frac{1}{\tau }\frac{z(\tau ,)}{|z(\tau ,)|}\rho (\tau ,)v_{C^0(I_ϵ,^2)}C\tau ^1$$

for $`\tau T`$. Then there is an $`\eta >0`$ and an isometry $`\varphi `$ from $`(D,g_D)`$ to $`(^2,g_R)`$ such that if $`(Q,P)=\varphi z`$, then the conclusions of Theorem 1 hold if we replace $`S^1`$ with $`I_\eta `$.

Proof. Using the isometries (8) and (19), we conclude that the relation between variables in the disc $`z`$ and the $`PQ`$-variables given by

(50)
$$(Q,P)=[\frac{2\mathrm{I}\mathrm{m}z}{1+|z|^22\mathrm{R}\mathrm{e}z},\mathrm{ln}(1|z|^2)+\mathrm{ln}(1+|z|^22\mathrm{R}\mathrm{e}z)]$$

is an isometry. Using (22) and (50), we have

(51)
$$P=\rho 2\mathrm{ln}(1+|z|)+\mathrm{ln}(1+|z|^22\mathrm{R}\mathrm{e}z).$$

By the assumptions of the lemma, there is a positive $`\gamma `$ such that

(52)
$$2\gamma |v(\theta )|12\gamma $$

for $`\theta I_ϵ`$ and by making $`ϵ`$ smaller, if necessary, we can assume that the image of $`I_ϵ`$ under $`v/|v|`$ is contained in an angle interval of length less that $`\pi `$. We can also assume that $`\rho 1`$ in $`[\tau _0,\mathrm{})\times I_ϵ`$. By (49) we have

$$\frac{z(\tau ,)}{|z(\tau ,)|}\frac{v}{|v|}_{C^0(I_ϵ,^2)}C\tau ^1$$

for $`\tau \tau _0`$. Consequently, we can assume that $`z/|z|`$ is contained in an angle interval of $`3\pi /2`$. Thus we can perform a rotation to the solution in order to ensure that

(53)
$$|z/|z|1|c$$

for some positive constant $`c`$ and $`(\tau ,\theta )[\tau _0,\mathrm{})\times I_ϵ`$. From now on, we will assume that such a rotation has been carried out. By (51) and (53), we conclude that $`P=\rho +O(1)`$. Combining this with (49) and (52), we get the conclusion that there is a $`T`$ such that for $`(\tau ,\theta )[T,\mathrm{})\times I_ϵ`$,

$$1P(\tau ,\theta )\tau 1,\gamma \frac{P(\tau ,\theta )}{\tau }1\gamma .$$

We can now carry out a localized version of the argument presented in the proof of Lemma 1. Instead of considering the set $`[T,\mathrm{})\times S^1`$ we consider $`(\tau ,\theta )`$ such that $`\tau T`$ and

$$\theta I_{ϵ,T,\tau }=(\theta _0ϵe^\tau +e^T,\theta _0+ϵ+e^\tau e^T),$$

where we assume that $`T`$ is big enough that $`e^Tϵ/2`$. We conclude that if

$$F_{ϵ,T}(\tau )=\frac{1}{2}e^\tau \underset{\theta I_{ϵ,T,\tau }}{sup}𝒜(\tau ,\theta )+\frac{1}{2}e^\tau \underset{\theta I_{ϵ,T,\tau }}{sup}(\tau ,\theta ),$$

where $`𝒜`$ and $``$ are defined in (24), then

$$F_{ϵ,T}(\tau )C\tau ^2.$$

For any $`\xi >0`$ there is a $`T_2`$ such that if $`\tau T_2`$, then we can change the initial data outside of $`(\theta _0ϵ/4,\theta _0+ϵ/4)`$ so that $`F(\tau )\xi `$, where $`F`$ is defined in (11) and $`\gamma P(\tau ,)/\tau 1\gamma `$. Let us be explicit concerning the change of the initial data. Let $`\chi C^{\mathrm{}}(S^1,)`$ be such that $`\chi (\theta )=1`$ for $`\theta `$ outside of $`(\theta _0ϵ/2,\theta _0+ϵ/2)`$ and $`\chi (\theta )=0`$ for $`\theta (\theta _0ϵ/4,\theta _0+ϵ/4)`$. Since $`Q`$ is bounded in the relevant set, due to (50) and (53), the modifications

$$\stackrel{~}{P}=(1\chi )P+\chi \gamma \tau ,\stackrel{~}{P}_\tau =(1\chi )P_\tau +\chi \gamma ,\stackrel{~}{Q}=(1\chi )Q,\stackrel{~}{Q}_\tau =(1\chi )Q_\tau $$

yield the desired conclusion for large enough $`\tau `$. By modifying the initial data as above at a late enough time $`T_3`$, the conditions of Theorem 1 will be fulfilled. Furthermore, we can assume that we have the original solution in the strip $`[T_3,\mathrm{})\times (\theta _0ϵ/8,\theta _0+ϵ/8)`$. Consequently, we get the conclusions of Theorem 1 localized to the above mentioned strip. $`\mathrm{}`$

###### Lemma 7.

Consider a solution to (20). Let $`\theta _0S^1`$ be a fixed angle, and assume that there is a $`T`$ and a $`0<\gamma <1`$ such that

$$\frac{\rho (\tau ,\theta _0)}{\tau }1\gamma $$

for all $`\tau T`$. Then there is an $`\eta >0`$ and a $`vC^0(I_\eta ,^2)`$ such that for $`\tau T`$,

$$\frac{1}{\tau }\frac{z(\tau ,)}{|z(\tau ,)|}\rho (\tau ,)v_{C^0(I_\eta ,^2)}+\frac{2z_\tau (\tau ,)}{1|z(\tau ,)|^2}v_{C^0(I_\eta ,^2)}$$

$$+e^\tau \frac{2z_\theta (\tau ,)}{1|z(\tau ,)|^2}_{C^0(I_\eta ,^2)}C\tau ^1.$$

Proof. Let

$$𝒮_\tau =(\theta _02e^\tau ,\theta _0+2e^\tau ),$$

where we assume that $`\tau `$ is big enough that $`2e^\tau \pi /2`$. We know that

$$|z_\tau |_D^2+e^{2\tau }|z_\theta |_D^2C$$

for $`(\tau ,\theta )[T,\mathrm{})\times S^1`$ and that (45) holds, and a similar division for the spatial derivatives, assuming $`|z|>0`$. Consequently,

$$|\frac{\rho (\tau ,\theta _1)}{\tau }\frac{\rho (\tau ,\theta _2)}{\tau }|\frac{C}{\tau },$$

assuming $`\theta _i𝒮_\tau `$. Thus there is a $`T_1T`$ such that (assuming $`3\beta =\gamma `$)

(54)
$$\rho (\tau ,\theta )/\tau 12\beta \mathrm{and}\rho (\tau ,\theta )\tau 2$$

for $`\tau [T_1,\mathrm{})`$ and $`\theta 𝒮_\tau `$. The idea of the argument is as follows. First we consider a quantity similar to the $`G`$ introduced in the proof of Lemma 5, the only difference being that we take supremum over $`𝒮_\tau `$ instead of $`S^1`$. At a late enough time, say $`\tau _1`$, the quantity analogous to $`G`$ has become small enough that we can close the argument and make statements concerning the domain determined by $`𝒮_{\tau _1}`$. This domain contains an open subset of the singularity.

Let $`_1`$ be defined as in (42) and similarly for $`_2`$. Let $`\tau \tau _0T_1`$. Define

$$L_1(u,\theta )=_1(u,\theta +e^ue^\tau )\mathrm{and}L_2(u,\theta )=_2(u,\theta e^u+e^\tau ),$$

where $`\theta 𝒮_\tau `$. Note that

$$\theta +e^ue^\tau ,\theta e^u+e^\tau 𝒮_u$$

for $`u\tau `$ and $`\theta 𝒮_\tau `$. Let

$$\widehat{L}_i(u)=\underset{\theta 𝒮_u}{sup}_i(u,\theta )\mathrm{and}\widehat{L}=\widehat{L}_1+\widehat{L}_2.$$

Due to (46) we have, for $`\theta 𝒮_\tau `$,

$$_1(\tau ,\theta )=L_1(\tau _0,\theta )+_{\tau _0}^\tau [(_ue^u_\theta )_1](u,\theta +e^ue^\tau )𝑑u$$

$$\widehat{L}_1(\tau _0)+_{\tau _0}^\tau (\frac{1}{2}\frac{1}{u})(\widehat{L}_1+\widehat{L}_2)𝑑u.$$

Taking the supremum over $`\theta 𝒮_\tau `$ and adding a similar estimate for $`_2`$, we get the conclusion that

$$\widehat{L}(\tau )\widehat{L}(\tau _0)+_{\tau _0}^\tau (1\frac{2}{u})\widehat{L}(u)𝑑u.$$

Consequently,

$$e^\tau \widehat{L}(\tau )e^{\tau _0}\widehat{L}(\tau _0)\frac{\tau _0^2}{\tau ^2}.$$

Let us consider a similar argument on a different set. Let

$$𝒮_{\tau _1,\tau }=(\theta _0e^\tau e^{\tau _1},\theta _0+e^\tau +e^{\tau _1}).$$

Note that $`𝒮_{\tau _1,\tau _1}=𝒮_{\tau _1}`$. For $`\theta 𝒮_{\tau _1,\tau }`$, let

$$K_1(u,\theta )=_1(u,\theta +e^ue^\tau )\mathrm{and}K_2(u,\theta )=_2(u,\theta e^u+e^\tau ).$$

For $`u\tau `$, we have

$$\theta +e^ue^\tau ,\theta e^u+e^\tau 𝒮_{\tau _1,u}.$$

Letting

$$\widehat{K}_i(u)=\underset{\theta 𝒮_{\tau _1,u}}{sup}_i(u,\theta )\mathrm{and}\widehat{K}=\widehat{K}_1+\widehat{K}_2,$$

we can argue similarly to the above in order to obtain

(55)
$$e^\tau \widehat{K}(\tau )e^{\tau _1}\widehat{K}(\tau _1)\frac{\tau _1^2}{\tau ^2},$$

assuming $`\rho (u,\theta )u2`$ for $`u[\tau _1,\tau ]`$ and $`\theta 𝒮_{\tau _1,u}`$. Note that $`\widehat{K}(\tau _1)=\widehat{F}(\tau _1)`$. The problem is of course that we cannot assume that $`\rho \tau 2`$ in the relevant set. Let us assume that $`\tau _1`$ is big enough so that

(56)
$$e^{\tau _1}\widehat{K}(\tau _1)\beta ^2,\tau _1\beta 2$$

and $`\tau _1T_1`$. Due to (54), we know that $`\rho (\tau _1,\theta )/\tau _112\beta `$ for $`\theta 𝒮_{\tau _1}`$. Let

$$𝒮=\{\tau [\tau _1,\mathrm{}):s[\tau _1,\tau ],\theta 𝒮_{\tau _1,s}\frac{\rho (s,\theta )}{s}1\beta \}.$$

Note that $`𝒮`$ is closed, connected and non-empty. Let us prove that it is open. Let $`\tau 𝒮`$. Then (55) is applicable in $`[\tau _1,\tau ]`$ due to (56). Consequently, if $`\theta 𝒮_{\tau _1,\tau }`$, then

$$|\frac{\rho (\tau ,\theta )}{\tau }\frac{\rho (\tau _1,\theta )}{\tau _1}|=|_{\tau _1}^\tau \frac{1}{s}(\rho _s\frac{\rho }{s})𝑑s|_{\tau _1}^\tau \beta \frac{\tau _1}{s^2}𝑑s=\beta (1\frac{\tau _1}{\tau }).$$

This implies, using (54), that

$$\frac{\rho (\tau ,\theta )}{\tau }1\beta (1+\frac{\tau _1}{\tau }).$$

Consequently, $`𝒮`$ is open. The conclusions of the lemma follow by an argument similar to the proof of Lemma 5. $`\mathrm{}`$

## Acknowledgements

This work was partly carried out when the author was enjoying the the hospitality of Rutgers. The author would like to express his gratitude to A. Shadi Tahvildar-Zadeh for the invitation and many stimulating discussions.


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# References

YITP-03-11

hep-th/0303240

Unitary-Matrix Integration on 2D Yang-Mills Action
Yoshinobu HABARA ***e-mail: habara@yukawa.kyoto-u.ac.jp

Yukawa Institute for Theoretical Physics,
Kyoto University, Kyoto 606-8502, Japan

## Abstract

Using the idea of Itzykson-Zuber integral, unitary-matrix integration of 2D Yang-Mills action is presented. The uniqueness of the solution of heat equation enables us to integrate out the unitary-matrix parts of hermite matrices and to obtain the expression of integration over vectors, also in this case.

Introduction

In 1980, C. Itzykson and J.-B. Zuber presented the method of integration over angular variables on a simplest model coupling two hermite matrices . Here, the angular variable means the unitary matrix which diagonalizes physical quantitiy, i.e. $`N\times N`$ hermite matrix. Their work had a great impact and enabled developments especially on the theories of random matrix and $`2D`$ gravity (  etc.). While, until then, there was Creutz’s unitary-matrix integration  well known in lattice gauge theory which integrate unitary matrices order by order in high temperature expansion, Itzykson-Zuber integral is epoch-making in the sense that they obtained the expression only using the eigenvalues of hermite matrices, after exactly integrating out angular variables, to all orders of exponential. In practice, since it is too complicated to integrate in usual ways, they evaluate the integral with the partial differential equation known as heat equation whose uniqueness of its solution is completely understood. They showed that the expression after integration is equivalent to that before integration using the uniqueness of the solution. After the unitary-matrix integration, the number of variables is reduced to $`N`$, i.e. vector, from $`N\times N`$ matrix, and therefore, the analyses such as $`N\mathrm{}`$ planar limit becomes very easy.

In this paper, using their idea, we present a similar expression after integrating out unitary matrices on 2-dimensional Yang-Mills action with intent to prepare to apply this, in future, to 10-dimensional Yang-Mills theory, i.e. IIB Matrix Model  which is known to be convergent . On our model, if we introduce an external field as $`N\times N`$ hermite matrix, the same heat equation as Itzykson-Zuber integral is also found to be satisfied, and therefore, we can evaluate our integral on the basis of the uniqueness of its solution. However, unlike Itzykson-Zuber integral, since a number of coefficients which cannot be determined only from the uniqueness remain, we derive them from the structure of the commutator characteristic of Yang-Mills action. As a result, in case $`N`$ is even, after integrating out angular variables, we find that the partition function of 2-dimensional Yang-Mills action can be expressed in terms of the remaining eivenvalues.

Unitary-Matrix Integration of 2D Yang-Mills Action

The integral we investigate is the Yang-Mills action in 2 “space-time” dimensions

$`I(C;g)={\displaystyle 𝑑A𝑑B\mathrm{exp}\left\{gtr[A,B]^2+itrC[A,B]\right\}},`$ (1)
which exists in the “external field” $`C`$. Here, $`A,B`$ and $`C`$ are $`N\times N`$ hermitian matrices, and $`iC[A,B]`$ is also hermitian. $`dA`$ (and $`dB`$) is the measure

$`dA={\displaystyle \underset{i=1}{\overset{N}{}}}dA_{ii}{\displaystyle \underset{i<j}{\overset{N}{}}}d(ReA_{ij})d(ImA_{ij}),`$ (2)
which is invariant under the adjoint operation of $`U(N)`$ matrices

$`d(V^{}AV)=dA,VU(N).`$
This measure can be written in the diagonalized form as follows:

$`dA={\displaystyle \underset{i=1}{\overset{N}{}}}da_idU\mathrm{\Delta }^2(\stackrel{}{a}),`$ (3)
where $`A`$ is diagonalized by unitary matrix $`U`$ such as $`A=U^{}aU,a=diag(a_1,\mathrm{},a_N)`$, and $`dU`$ is the normalized $`U(N)`$ Haar measure, $`𝑑U=1`$. $`\mathrm{\Delta }(\stackrel{}{a})`$ is the Vandermonde determinant

$`\mathrm{\Delta }(\stackrel{}{a})=(1)^{\frac{N(N1)}{2}}\left|\begin{array}{cccc}1& 1& \mathrm{}& 1\\ a_1& a_2& \mathrm{}& a_N\\ a_1^2& a_2^2& \mathrm{}& a_N^2\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\\ a_1^{N1}& a_2^{N1}& \mathrm{}& a_N^{N1}\end{array}\right|={\displaystyle \underset{i<j}{\overset{N}{}}}(a_ia_j).`$ (9)
As we can check easily, the Itzykson-Zuber integral ,

$`K(C;g)`$ $`=`$ $`{\displaystyle 𝑑X\mathrm{exp}\left\{gtrX^2+trCX\right\}}`$ (10)
$`=`$ $`{\displaystyle d^N\stackrel{}{x}\frac{\mathrm{\Delta }(\stackrel{}{x})}{\mathrm{\Delta }(\stackrel{}{c})}\mathrm{exp}\left\{g|\stackrel{}{x}|^2+\stackrel{}{c}\stackrel{}{x}\right\}},`$
obeys the following partial differenial equation,

$`tr\left({\displaystyle \frac{^2}{C^2}}\right)K(C;g)={\displaystyle \frac{}{g}}K(C;g),`$ (11)
$`tr\left({\displaystyle \frac{^2}{C^2}}\right){\displaystyle \underset{i,j=1}{\overset{N}{}}}{\displaystyle \frac{}{C_{ij}}}{\displaystyle \frac{}{C_{ji}}}.`$
Essentially, this equation is the heat equation, and has an unique solution (5) under the boundary condition $`K(C;g\mathrm{})=0`$.

Similarly, the integral (1) satisfies the same equation

$`tr\left({\displaystyle \frac{^2}{C^2}}\right)I(C;g)={\displaystyle \frac{}{g}}I(C;g),`$ (12)
and must have an unique solution under the boundary condition $`I(C;g\mathrm{})=0`$. After integrating over $`U(N)`$ matrices, we expect that (1) becomes the integral which only contains the integration over two eigenvalues of hermite matrices $`A`$ and $`B`$. Then, we suppose the integral (1) to become like the following form:

$`i(\stackrel{}{c};g)={\displaystyle d^N\stackrel{}{a}d^N\stackrel{}{b}\mathrm{exp}\left\{g\underset{i,j,k,l=1}{\overset{N}{}}R_{ijkl}a_ib_ja_kb_l+\underset{i,j,k=1}{\overset{N}{}}\mathrm{\Gamma }_{ijk}c_ia_jb_k\right\}}.`$ (13)
Here, the coefficients $`R_{ijkl}`$ and $`\mathrm{\Gamma }_{ijk}`$ are independent of $`\stackrel{}{c}=(c_1,\mathrm{},c_N)`$. Since this integral must obey the partial differential equation

$`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{^2}{c_j^2}}i(\stackrel{}{c};g)={\displaystyle \frac{}{g}}i(\stackrel{}{c};g),`$ (14)
$`R_{ijkl}`$ is determined from $`\mathrm{\Gamma }_{ijk}`$:

$`R_{ijkl}={\displaystyle \underset{m=1}{\overset{N}{}}}\mathrm{\Gamma }_{mij}\mathrm{\Gamma }_{mkl}.`$ (15)
Therefore, it is sufficient we obtain $`\mathrm{\Gamma }_{ijk}`$, and, from the equation

$`tr\left({\displaystyle \frac{^2}{C^2}}\right)=\mathrm{\Delta }^1(\stackrel{}{c}){\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{^2}{c_j^2}}\mathrm{\Delta }(\stackrel{}{c})+_{U_C}^2,`$ (16)
$`C=U_C^{}cU_C,`$
where $`_{U_C}^2`$ only depends on $`U_C`$, the integral (1) becomes

$`I(C;g)={\displaystyle d^N\stackrel{}{a}d^N\stackrel{}{b}\frac{f(\stackrel{}{a},\stackrel{}{b})}{\mathrm{\Delta }(\stackrel{}{c})}\mathrm{exp}\left\{g\underset{i,j,k,l=1}{\overset{N}{}}R_{ijkl}a_ib_ja_kb_l+\underset{i,j,k=1}{\overset{N}{}}\mathrm{\Gamma }_{ijk}c_ia_jb_k\right\}},`$ (17)
owing to the uniqueness of the solution of (7). Here, $`f(\stackrel{}{a},\stackrel{}{b})`$ is the undetermined function of $`\stackrel{}{a}`$ and $`\stackrel{}{b}`$. Furthermore, though $`\stackrel{}{c}`$ is exactly the eigenvalue of hermite matrix $`C`$, we cannot certain that both $`\stackrel{}{a}`$ and $`\stackrel{}{b}`$ are those of $`A`$ and $`B`$, yet. However, with $`g=0`$ in (12), the integral

$`{\displaystyle 𝑑U𝑑V}`$ $`\mathrm{exp}\left\{itrC[U^{}AU,V^{}BV]\right\}`$ (18)
$`={\displaystyle 𝑑U𝑑V𝑑W\mathrm{exp}\left\{itrW^{}CW[U^{}AU,V^{}BV]\right\}}`$
$`={\displaystyle 𝑑V𝑑W\mathrm{exp}\left\{itrA[V^{}BV,W^{}CW]\right\}}`$
$`={\displaystyle 𝑑U𝑑W\mathrm{exp}\left\{itrB[W^{}CW,U^{}AU]\right\}}`$
$`={\displaystyle \frac{1}{\mathrm{\Delta }(\stackrel{}{a})\mathrm{\Delta }(\stackrel{}{b})\mathrm{\Delta }(\stackrel{}{c})}}\mathrm{exp}\left\{{\displaystyle \underset{i,j,k=1}{\overset{N}{}}}\mathrm{\Gamma }_{ijk}c_ia_jb_k\right\}`$
enables us to derive $`f(\stackrel{}{a},\stackrel{}{b})=\mathrm{\Delta }(\stackrel{}{a})\mathrm{\Delta }(\stackrel{}{b})`$ and make sure that both $`\stackrel{}{a}`$ and $`\stackrel{}{b}`$ are the eigenvalues of $`A`$ and $`B`$, because all $`A,B`$ and $`C`$ stand on the same ground. And, also from the facts that $`trC[A,B]=trA[B,C]=trB[C,A]`$ and $`trC[A,B]=trC[B,A]`$, the index of $`\mathrm{\Gamma }_{ijk}`$ has cyclic symmetry: $`\mathrm{\Gamma }_{ijk}=\mathrm{\Gamma }_{jki}=\mathrm{\Gamma }_{kij}`$ and anti-symmetry under the interchange of its indices: $`\mathrm{\Gamma }_{ijk}=\mathrm{\Gamma }_{ikj}`$. Therefore, we obtain

$`I(C;g)`$ $`={\displaystyle 𝑑A𝑑B\mathrm{exp}\left\{gtr[A,B]^2+itrC[A,B]\right\}}`$ (19)
$`={\displaystyle d^N\stackrel{}{a}d^N\stackrel{}{b}\frac{\mathrm{\Delta }(\stackrel{}{a})\mathrm{\Delta }(\stackrel{}{b})}{\mathrm{\Delta }(\stackrel{}{c})}\mathrm{exp}\left\{g\underset{i,j,k,l=1}{\overset{N}{}}R_{ijkl}a_ib_ja_kb_l+\underset{i,j,k=1}{\overset{N}{}}\mathrm{\Gamma }_{ijk}c_ia_jb_k\right\}}.`$
The remaining task we must do is to determine the coefficients $`\mathrm{\Gamma }_{ijk}`$. Since to have the matrices $`A`$ and $`B`$ multiplied by the constants $`k`$ is equivalent to make the eigenvalues $`k\stackrel{}{a}`$ and $`k\stackrel{}{b}`$ respectively, $`\mathrm{\Gamma }_{ijk}`$ are also independent of $`\stackrel{}{a}`$ and $`\stackrel{}{b}`$. Then, we can evaluate $`\mathrm{\Gamma }_{ijk}`$ through the structure peculiar to the commutator $`[,]`$. We can easily check the integral

$`I(\stackrel{}{a},\stackrel{}{b};g)`$ $`={\displaystyle 𝑑U_A𝑑U_B\mathrm{exp}\left\{gtr[U_A^{}aU_A,U_B^{}bU_B]^2\right\}}`$ (21)
$`=\mathrm{exp}\left\{g{\displaystyle \underset{i,j,k,l=1}{\overset{N}{}}}R_{ijkl}a_ib_ja_kb_l\right\}`$
follows from (14), and (15) satisfies the differential equation due to the characteristic structure of the commutator:

$`{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{}{a_i}}I(\stackrel{}{a},\stackrel{}{b};g)=0,`$ (22)
therefore, from (16), we obtain the algebraic equation which $`R_{ijkl}`$ must obey:

$`{\displaystyle \underset{i,j,k,l=1}{\overset{N}{}}}R_{ijkl}b_ja_kb_l={\displaystyle \underset{i,j,k,l,m=1}{\overset{N}{}}}\mathrm{\Gamma }_{mij}\mathrm{\Gamma }_{mkl}b_ja_kb_l=0.`$ (23)
This is the identity about $`a_i`$ and $`b_i`$. By introducing $`d_m_{k,l}\mathrm{\Gamma }_{mkl}a_kb_l`$, this reduces to the identity

$`{\displaystyle \underset{i,j,m=1}{\overset{N}{}}}\mathrm{\Gamma }_{mij}b_jd_m=0`$ (24)
about $`b_j`$ and $`d_m`$, and we obtain conditions for $`\mathrm{\Gamma }_{ijk}`$:

$`{\displaystyle \underset{i=1}{\overset{N}{}}}\mathrm{\Gamma }_{ijk}=0,\text{for fixed }j,k.`$ (25)
When $`N=4`$, these equations become

$`\mathrm{\Gamma }_{312}+\mathrm{\Gamma }_{412}=0,\mathrm{\Gamma }_{123}+\mathrm{\Gamma }_{423}=0,`$
$`\mathrm{\Gamma }_{213}+\mathrm{\Gamma }_{413}=0,\mathrm{\Gamma }_{124}+\mathrm{\Gamma }_{324}=0,`$
$`\mathrm{\Gamma }_{214}+\mathrm{\Gamma }_{314}=0,\mathrm{\Gamma }_{134}+\mathrm{\Gamma }_{234}=0,`$
then, except overall factor, all $`\mathrm{\Gamma }_{ijk}`$s are determined exactly as

$`\mathrm{\Gamma }_{123}=\mathrm{\Gamma }_{124}=\mathrm{\Gamma }_{134}=\mathrm{\Gamma }_{234}=1.`$ (26)
When $`N=2`$ and $`N=3`$, $`\mathrm{\Gamma }_{ijk}=0`$ for all $`i,j,k`$. When $`N4`$, from the symmetry of its indices, $`\mathrm{\Gamma }_{ijk}`$ has $`{}_{N}{}^{}C_{3}^{}`$ independent components, but (20) are $`{}_{N}{}^{}C_{2}^{}`$ equations at most. So, we cannot determine all $`\mathrm{\Gamma }_{ijk}`$ from (20) in general. However, as is obvious from (14), all $`a_i,b_i`$ and $`c_i`$ are contained in the same weight. That is, except the overall factor which can be absorbed into the integration measure, we can set $`\left|\mathrm{\Gamma }_{ijk}\right|=1`$. Then, for even $`N4`$, we can find the solution of (20):

$`\mathrm{\Gamma }_{ijk}=(1)^{i+j+k},\text{for}i<j<k.`$ (27)
When $`N`$ is odd, the solution such as $`\left|\mathrm{\Gamma }_{ijk}\right|=1`$ does not exist, because odd numbers of $`\pm 1`$ cannot add up to zero. In fact, when $`N=5`$,

$`\mathrm{\Gamma }_{312}+\mathrm{\Gamma }_{412}+\mathrm{\Gamma }_{512}=0,\mathrm{\Gamma }_{124}+\mathrm{\Gamma }_{324}+\mathrm{\Gamma }_{524}=0,`$
$`\mathrm{\Gamma }_{213}+\mathrm{\Gamma }_{413}+\mathrm{\Gamma }_{513}=0,\mathrm{\Gamma }_{125}+\mathrm{\Gamma }_{325}+\mathrm{\Gamma }_{425}=0,`$
$`\mathrm{\Gamma }_{214}+\mathrm{\Gamma }_{314}+\mathrm{\Gamma }_{514}=0,\mathrm{\Gamma }_{134}+\mathrm{\Gamma }_{234}+\mathrm{\Gamma }_{534}=0,`$
$`\mathrm{\Gamma }_{215}+\mathrm{\Gamma }_{315}+\mathrm{\Gamma }_{415}=0,\mathrm{\Gamma }_{135}+\mathrm{\Gamma }_{235}+\mathrm{\Gamma }_{435}=0,`$
$`\mathrm{\Gamma }_{123}+\mathrm{\Gamma }_{423}+\mathrm{\Gamma }_{523}=0,\mathrm{\Gamma }_{145}+\mathrm{\Gamma }_{245}+\mathrm{\Gamma }_{345}=0,`$
the solution does not exist. This strange feature that we can use the expression (14) only for even $`N`$ seems to come from the following characteristic of the commutator $`[,]`$. The algebra of hermite matrices has the subalgebra which consists of real symmetric matrices $`A^{}=A^T=A`$, and, in this subalgebra, the commutator $`[A,B]`$ always becomes the anti-symmetric matrix $`[A,B]^T=[A,B]`$ which has even rank. That is, the commutator $`[,]`$ essentially contains the structure of even rank which is hidden by larger $`U(N)`$ symmetry before the integration, and then, the integration of $`U(N)`$ matrices may reveal this structure. However, since we need the integral after taking the limit $`N\mathrm{}`$, it is sufficient for our purpose to have the result that $`N`$ is even number.

Lastly, let us line up the result we obtain in this paper.

$`I(C;g)`$ $`={\displaystyle 𝑑A𝑑B\mathrm{exp}\left\{gtr[A,B]^2+itrC[A,B]\right\}}`$
$`={\displaystyle d^N\stackrel{}{a}d^N\stackrel{}{b}\frac{\mathrm{\Delta }(\stackrel{}{a})\mathrm{\Delta }(\stackrel{}{b})}{\mathrm{\Delta }(\stackrel{}{c})}\mathrm{exp}\left\{g\underset{i,j,k,l=1}{\overset{N}{}}R_{ijkl}a_ib_ja_kb_l+\underset{i,j,k=1}{\overset{N}{}}\mathrm{\Gamma }_{ijk}c_ia_jb_k\right\}},`$
$`R_{ijkl}`$ $`={\displaystyle \underset{m=1}{\overset{N}{}}}\mathrm{\Gamma }_{mij}\mathrm{\Gamma }_{mkl},`$
$`\mathrm{\Gamma }_{ijk}`$ $`=\mathrm{\Gamma }_{kij}=\mathrm{\Gamma }_{ikj},`$
$`\mathrm{\Gamma }_{ijk}`$ $`=(1)^{i+j+k},\text{for}i<j<k.`$
Discussions

Using the uniqueness of the solution for heat equation which plays a crucial role on Itzykson-Zuber integral, we derive an expression which contains eigenvalues of hermite matrices only as integration variables after integrationg out angular variables on 2-dimensional Yang-Mills action. Then, characteristic structures of the commutator are important and needed to determine coefficients in the expression. And the expression is possible only for even $`N`$, reflecting the structures.

Our expression can be applied to various problems. For example, as the model of Itzykson-Zuber integral which can be solved exactly in $`N\mathrm{}`$ planar limit, we can evaluate similarly on 2-dimensional Yang-Mills model. And, as refered in introduction, by extending our result to larger dimensions, for example IIB Matrix Model, we may be able to treat the phenomenon of the phase-transition of space-time, i.e. inflation. However, when we extend to higher dimensions than three, we cannot probably use the uniqueness of the solution of heat equation, so must introduce other methods.

Furthermore, since Vandermonde determinant $`\mathrm{\Delta }(\stackrel{}{x})`$, which always emerges in case of integrating over hermite matrices, is the wave function of ground state of Calogero-Sutherland model describing the system of fermions on circle , and the orthogonal polynomials containing $`\mathrm{\Delta }(\stackrel{}{x})`$ satisfy particular recursion relations , we may be able to calculate exactly over remaining eigenvalues.

Acknowledgements

I would like to thank my colleagues for useful comments and discussions. Especially, I gratefully acknowledge S. Matsuura and S. Shinohara for helpful suggestions.


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# CONTRASTING THE UV AND X-RAY O vi COLUMN DENSITY INFERRED FOR THE OUTFLOW IN NGC 5548

## 1 INTRODUCTION

Outflows in Seyfert galaxies are evident by resonance line absorption troughs, which are blueshifted with respect to the systemic redshift of their emission counterparts. Velocities of several hundred km s-1 (Crenshaw et al. 1999; Kriss et al. 2000) are observed in both UV resonance lines (e.g., C iv $`\lambda \lambda `$1548.20,1550.77, N v $`\lambda \lambda `$1238.82,1242.80, O vi $`\lambda \lambda `$1031.93,1037.62 and Ly$`\alpha `$), as well as in X-ray resonance lines (Kaastra et al. 2000; Kaspi et al. 2000). Similar outflows (often with significantly higher velocities) are seen in quasars which are the luminous relatives of Seyfert galaxies (Weymann et al. 1991; Korista, Voit, Morris, & Weymann 1993; Arav et al. 2001a). Reliable measurement of the absorption column densities in the troughs are crucial for determining the ionization equilibrium and abundances of the outflows, and the relationship between the UV and the ionized X-ray absorbers.

NGC 5548 is one of the most studied Seyfert galaxies, including intensive reverberation campaigns (Netzer & Maoz 1990; Clavel et al. 1991; Peterson et al. 1991; Korista et al. 1995), line studies (Krolik et al. 1991; Ferland et al. 1992; Goad & Koratkar 1998; Kaspi & Netzer 1999; Korista & Goad 2000), and theoretical modeling (Done & Krolik 1996; Bottorff, Korista & Shlosman 2000; Srianand 2000). The intrinsic absorber in NGC 5548 was studied in the UV using the International Ultraviolet Explorer (Shull & Sachs 1993), the HST Goddard High Resolution Spectrograph (GHRS) (Mathur, Elvis & Wilkes 1999) and Space Telescope Imaging Spectrograph (STIS) (Crenshaw & Kraemer 1999), the Far Ultraviolet Spectroscopic Explorer (FUSE) (Brotherton et al. 2002), and in X-ray with the ASCA (George et al. 1998) and Chandra (Kaastra et al. 2000, 2002) satellites. These high quality observations combined with the relative simplicity of its intrinsic absorption features make NGC 5548 an excellent target for understanding the nature of Seyfert outflows.

In the last few years our group (Arav 1997; Arav et al. 1999a; Arav et al. 1999b; de Kool et al. 2001; Arav et al. 2001) and others (Barlow 1997, Telfer et al. 1998, Churchill et al. 1999, Ganguly et al. 1999) have shown that in quasar outflows most lines are saturated even when not black. We have also shown that in many cases the shapes of the troughs are almost entirely due to changes in the line of sight covering as a function of velocity, rather than to differences in optical depth (Arav et al. 1999b; de Kool et al. 2001; Arav et al. 2001a). Gabel et al. (2002) show the same effect in the outflow troughs of NGC 3783. As a consequence, the column densities inferred from the depths of the troughs are only lower limits.

These results have led us to suspect that the current determination of column densities in Seyfert outflows is highly uncertain. Recently, Arav, Korista & de Kool (2002) have re-analyzed the HST high resolution spectroscopic data of the intrinsic absorber in NGC 5548 and found that the C iv absorption column density in the main trough is at least four times larger than previously determined. Furthermore, in the same paper it was shown that similar to the case in quasars, the shape of the main trough is almost entirely due to changes in covering fraction as a function of velocity, and is not due to varying column density as a function of velocity.

An important new saturation diagnostic is available in the form of UV and X-ray absorption troughs that arise from the same ion. The relationship between the ionic column density ($`\mathrm{N}_{\mathrm{ion}}`$) and the optical depth ($`\tau `$) of an absorption trough is given by (see eq. (3) in Arav et al. 1999a):

$$\mathrm{N}_{\mathrm{ion}}=\frac{3.8\times 10^{14}\mathrm{cm}^2}{\lambda _\mathrm{o}\mathrm{f}_{\mathrm{ik}}}\tau (\mathrm{v})\mathrm{dv},$$
(1)
where $`\lambda _\mathrm{o}`$ and $`f_{ik}`$ are the wavelength and oscillator strength of the transition, respectively, and $`v`$ is the velocity in km s-1. Using the apparent optical depth ($`I=e^\tau `$, where $`I`$ is the residual intensity) in equation (1), a trough with the same $`I(v)`$ yields fifty times higher $`\mathrm{N}_{\mathrm{ion}}`$ for an X-ray line as compared to a UV line with the same $`f_{ik}`$, simply because its wavelength is fifty times smaller. In cases where the UV trough is saturated and its shape is determined by changes in covering fraction as a function of velocity, it is virtually impossible to distinguish between $`\tau `$ values above $`3`$ (see Arav, Korista & de Kool 2002). For equal oscillator strength, the X-ray trough of the same ion will have $`\tau =1`$ when the UV trough reaches $`\tau _{UV}=50`$. Therefore, the X-ray trough is sensitive to a much higher column density than the UV one and can provide a crucial saturation diagnostic for the UV troughs.

In this paper we compare the O vi column density inferred independently from UV and X-ray absorption features in the spectrum of NGC 5548, and demonstrate the large discrepancy between the two values. The X-ray derived O vi column density is seven times larger than the UV one. Since the X-ray determination is a robust lower limit, we conclude that the UV lines are much more saturated than originally inferred, and identify the UV analysis as the probable cause for this discrepancy. Specifically, Gaussian modeling cannot account for a velocity dependent covering fraction, and more generally Gaussian modeling is a poor representation for absorption associated with a dynamical outflow.

## 2 CONTRASTING THE UV AND X-RAY INFERRED O vi COLUMN DENSITY

### 2.1 Data Description

NGC 5548 was observed by FUSE on 2000 June 7 with a total exposure time of 25 ksec. A full description of the observations, reduction process and plots of the spectra are presented in Brotherton et al. (2002). These observations include the spectral region associated with both the emission and absorption features of the O vi $`\lambda \lambda `$1032,1038 doublet in the rest frame of NGC 5548.

Our XMM-Newton data were taken on 2001 July 12. The total exposure was 116 ksec and the data were reduced using the standard SAS analysis package. Here we only discuss the Reflection Grating Spectrometer (RGS) spectra. A detailed description of the RGS instruments is given by Den Herder et al. (2001). Briefly, there are two almost identical RGS detectors aboard XMM-Newton that cover the 5$``$38 Å range with an almost uniform spectral resolution of 0.05 Å, and have a maximum effective area of about 140 cm2 at 15 Å. In the present work we focus upon the O vi line region around 22 Å. Due to failure of one of the 9 CCD’s of RGS2, only data from RGS1 are available in this wavelength band and analyzed here. The spectrum was fitted using a continuum plus absorption components as described by Kaastra et al. (2002). A full analysis of the spectrum will be given elsewhere (Steenbrugge et al., 2002). as well as timing analysis of the data (Markowitz et al. 2002) and analysis of the European Photon Imaging Camera (EPIC) data (Pounds et al. 2003).

### 2.2 O vi Column Density Extractions

In the XMM-Newton/RGS spectrum there is a distinct absorption line visible at a rest wavelength of 22.0 Å due to O vi (see Fig. 1). The line occurs in the wavelength range of the well-known He-like triplet of O vii, from which the forbidden line at 22.1 Å and the intercombination line at 21.8 Å are seen in emission, and the resonance line at 21.6 Å in absorption (cf. Kaastra et al. 2002). The O vi absorption line at 22.0 Å is in fact a satellite to these He-like lines, and its importance for AGN spectra was first recognized by Pradhan (2001). The line (in fact an unresolved blend of the so-called $`r`$ and $`u`$ lines) is caused by inner-shell photo-excitation of an O vi ion in the ground state to the 1s2p(1P)2s 2P1/2 and 1s2p(1P)2s 4P3/2 levels, respectively. We used atomic data for this transition as calculated by Behar (2002, in preparation). The wavelength for the blend of the two components is 22.007 Å and the combined oscillator strength $`f`$ is 0.525 (Pradhan gives 22.05 Å and $`f=0.576`$).

The equivalent width of this O vi trough is 40 $`\pm 10`$ mA, and therefore a lower limit for the associated column density is $`3.2\pm 0.8\times 10^{16}`$ cm-2, where the error is directly associated with the signal-to-noise of the data. A detailed description of the fitting method used to extract this O vi column density value is given in Steenbrugge et al. (2002)

We stress that this is a lower limit since: a) There may be unresolved absorption components due to the $`R=500`$ resolution of the RGS at this wavelength. b) The trough may show non-black saturation even for resolved components, similar to the situation evident in the UV troughs associated with the outflow. c) The adjacent O vii f emission line, can partly fill in the O vi trough. We note that the FWHM of the O vi trough is 800$`\pm `$200 km s-1, and that it’s centroid velocity is $`370\pm 160`$ km s-1 for a systemic redshift of $`z=0.017175`$ (see redshift discussion in Brotherton et al, 2002).

Brotherton et al. (2002) analyzed the FUSE spectra of NGC 5548 in order to determine the O vi absorption column density from the O vi $`\lambda \lambda `$1031.93,1037.62 (oscillator strength of 0.13 and 0.066, respectively; Verner, Verner & Ferland 1996). They used the program SPECFIT (Kriss 1994) to fit the O vi spectral region using the following assumptions: Approximating the intrinsic absorption as several components with Gaussian distributions in optical depth that are permitted to partially cover the background source, where the optical depth ratio of the doublets is constrained to be 2:1 at all velocities. They have looked at two types of models. one where the outflow covers all the emission components (continuum, broad emission lines, and narrow emission lines \[NELs\]), and one where the flow does not cover the NELs. Since the model with un-covered NELs gives higher $`N_{O\mathrm{vi}}`$ estimates, we will concentrate on it, noting that our discussion is similarly applicable to the other model.

The overall quality of the Brotherton et al. (2002) fit is very good (see their Fig. 3) and the statistical error quoted for the $`N_{O\mathrm{vi}}`$ associated with each of the absorption components is mostly between 10–20% (their Table 2). We note though, that their model includes 43 free parameters: Eight absorption Gaussian with four free parameters each (position, depth, width and covering fraction), three emission Gaussians with three free parameters each (position, hight and width), and two free parameters for the continuum (flux and slope). Although the goodness of the fit is impressive, it is not surprising given the number of free parameters involved. Summing up the $`N_{O\mathrm{vi}}`$ determinations in each component we obtain a total $`N_{O\mathrm{vi}}=4.9\pm 0.6\times 10^{15}`$ cm-2, where for a conservative error estimate we simply added the quoted individual errors. The total width of the O vi UV absorption complex is $`1200`$ km s-1, which is consistent with the width of the O vi X-ray trough (see above).

We therefore possess two very different estimates for the O vi absorption column density associated with the outflow in NGC 5548. On the one hand the XMM/RGS observation gives a lower limit $`N_{O\mathrm{vi}}3.2\pm 0.8\times 10^{16}`$ cm-2, on the other hand a Gaussian fitting of the FUSE data yields $`N_{O\mathrm{vi}}=4.9\pm 0.6\times 10^{15}`$ cm-2.

### 2.3 Kinematic Distribution of the X-ray-Infered O vi Column Density

At 22 Å the resolution element of the RGS instrument is roughly 700 km s-1. In order to obtain the deep O vi absorption trough seen in the data (Fig. 1), the absorbing material must therefore be spread over a similar velocity interval. An absorption component with a much narrower velocity width cannot account for the X-ray trough even if its optical depth is very high. This narrow yet deep intrinsic trough will be transformed into a shallow trough of roughly 700 km s-1 width by the line spread function of the instrument.

Brotherton et al. (2002) measured FWHM$`<150`$ km s-1 for UV components 2-6 of the outflow. Therefore, the high X-ray-inferred $`N_{O\mathrm{vi}}`$ cannot be hidden in only one, two, or three of these components even if they are highly saturated. UV components 0.5 and 1 are wider (350 and 280 km s-1, respectively). However, even combined, these components cannot by themselves account for the X-ray-inferred $`N_{O\mathrm{vi}}`$. Their combined width and shape are not enough to produce the observed X-ray trough even if we assume both components are highly saturated. We conclude that since the X-ray-inferred $`N_{O\mathrm{vi}}`$ must be associated with the UV O vi kinematic components, at least half of these components have to be highly saturated to allow for the observed X-ray O vi absorption trough.

## 3 WHAT CAUSES THE DISCREPANCY IN THE O vi COLUMN DENSITY DETERMINATIONS?

### 3.1 The Plausibility of Real Temporal Changes

A difference of 13 months exist between the FUSE (2000 June 7) and XMM (2001 July 12) observations of NGC 5548. It is therefore important to address the possibility that the order of magnitude difference in inferred $`N_{O\mathrm{vi}}`$ is caused by temporal changes. In table 1 we list the relevant UV and X-ray observations, including the continuum flux associated with each. Based on the four available epochs of high-resolution UV spectroscopy of the object, we conclude that strong variations in the shape and strength of the UV troughs over a year time-scale are unlikely. Components 2, 4 and 5 of the C iv trough (following Crenshaw & Kraemer 1999) show only marginal changes in the three available epochs of high-resolution spectroscopy: HST/GHRS spectrum taken in 1996 August 24 (Mathur, Elvis & Wilkes 1999), HST/STIS observation from 1998 March 11 (Crenshaw & Kraemer 1999) and HST/STIS observation from 2002 January 22 (Crenshaw et al. in preparation). Only components 1 and 3 (which are blended to some extent) show significant variation (at the $``$50% level). The main C iv trough (component 4) showed less than 10% changes in residual intensity between all three epochs. Furthermore, component 4 in the C iv trough is almost identical in shape and depth to the main component in the O vi trough covered by the FUSE observations (see Fig. 4 in Arav et al. 2002). Such similarity argues for both small temporal variation in the troughs, as well as for a moderate to high degree of saturation in both the C iv and O vi troughs (see discussion in Arav et al. 2002). These comparisons suggest that the shape of the main UV troughs in NGC 5548 change rather moderately on a time scale of 2 – 6 years.

We note that if the UV troughs show non-black saturation, large column density changes are possible even when the absorption troughs do not change appreciably. However, UV analysis techniques (Gaussian modeling or otherwise) have a limited dynamical range for inferring real optical depth (see § 3.3). Therefore, if the troughs’ shape did not change appreciably it is very difficult to deduce a real change in the O vi absorbing column density even if the level of saturation increased considerably. Furthermore, the X-ray data suggest that dramatic $`N_{O\mathrm{vi}}`$ changes are unlikely over a one year time scale.

Recently (January 2002) a very long Chandra LETGS observation (340 ksec) of NGC 5548 was obtained (Kaastra et al. in preparation). Preliminary analysis indicates that the detected O vi trough is very similar to the one seen in the XMM/RGS observation. Such occurrence suggests that an order of magnitude changes in the EW of the O vi X-ray trough between the FUSE (2000 June 7) and XMM (2001 July 12) observations of NGC 5548 is unlikely. There is also an earlier epoch Chandra LETGS observation taken on December 1999 (86 ksec; Kaastra et al. 2000) and have now been reanalyzed using improved wavelength and effective area calibration by Kaastra et al. (2002). The S/N of this observation is significantly lower than the S/N of either the XMM data (shown here), or the new Chandra LETGS observation. The O vi line is marginally detected (Fig. 3 in Kaastra et al. 2000), with a 1$`\sigma `$ upper limit of $`3.4\times 10^{16}`$ cm-2, which is consistent with the XMM measurement.

Additional support for relatively small changes in the O vi X-ray trough on a year timescale comes from comparing total fluxes in the 2-10 keV bands. The X-ray flux dropped by roughly 50% between the two higher quality observations (XMM and January 2002 Chandra), yet the derived $`N_{O\mathrm{vi}}`$ values are similar in both observations. This occurence suggest that the $`N_{O\mathrm{vi}}`$ value is quite insensitive to considerable flux changes over a year time scale.

### 3.2 Validity of the Underlying Physical Model

The modeling of the absorption troughs presented in Brotherton et al. (2002) is the most sophisticated and detailed Gaussian fitting of an AGN outflow to date. As mentioned above, the “goodness” of their fit is impressive. It is clear that a Gaussian fitting algorithm can obtain a low $`\chi ^2`$ fit between the data and a model spectrum. However, after obtaining a good fit the important question is how physically valid is the Gaussian fitting process? In other words, how well would the absorption column densities determined from Gaussian components represent the actual column densities of the AGN outflow?

Gaussian fitting is a common analysis tool for spectral absorption features. In particularly it is used in ISM and IGM studies. In these environments a Gaussian distribution of optical depth as a function of velocity seems physically plausible. This is because there is little dynamics involved, and therefore a thermal velocity distribution function yields a Gaussian optical depth function (even if “turbulent” broadening is invoked to explain the larger-than-thermal width). Furthermore, for ISM and IGM the geometrical size of the absorption clouds are much larger than the emission sources they cover. Therefore, the assumption that along our line of sight the absorber covers the emission source completely and uniformly is very reasonable.

When these models are adapted to represent AGN outflows two ingredients are added: a constant covering fraction to allow for partial line-of-sight-covering of the emission source, and a Gaussian velocity width much larger than is found for ISM and IGM clouds (in order to account for the width of the outflow’s absorption features). The physical applicability of both ingredients is highly questionable: First, careful analysis of high-resolution, high S/N spectra of AGN outflows reveals that the covering fraction is a strong function of velocity (Arav et al. 1999b; de Kool et al. 2001, Gabel et al. 2002). In particular, this is also the situation we find for the C iv trough of NGC 5548 (Arav et al. 2002). The strong dependence of the covering fraction upon velocity naturally undermines the constant covering fraction assumption.

Similarly, the large Gaussian velocity width invoked to fit the outflows features is also problematic. For acceptable photoionization equilibria, clouds’ thermal widths are $`<20`$ km s-1. For IGM clouds a “turbulent” broadening is invoked to explain the larger-than-thermal width. However, the FWHM of this broadening is rather small (typical values for the Doppler parameter are $`b`$ 16 – 46 km s-1; Sembach et al. 2001) and is associated with motions on scales of hundreds of kpc. In contrast the FWHM of Gaussian absorption components in the Brotherton et al. (2002) modeling reaches 350 km s-1 and is associated with length scales of several parsecs at most. The larger FWHM and the much smaller length scales associated with the AGN outflow weakens the physical plausibility of Gaussian optical depth distribution as compared with the IGM case. We conclude that there is little physical basis for modeling outflow troughs as Gaussians clouds with a constant covering fraction.

Both the strong dependence of the covering fraction upon velocity and the large width of the outflow absorption features can be simply accounted for by appealing to the dynamical nature of the outflow. Whereas it is reasonable to expect that a cloud will have a single covering factor (although, see de Kool et al. 2002), this is not the situation for an outflow, where the covering factor can be a strong function of velocity. Toy models which illustrate these kinematic/geometric effects are found in Arav (1996) and Arav et al. (1999a). A similar picture arises in global models for the structure of AGN (Elvis 2000; Elvis 2001).

### 3.3 Asymmetric Errors Associated with the Extraction of Column Densities

Another important issue is the estimate of errors associated with the extraction of column densities from the absorption troughs. A formal statistical error from minimizing the $`\chi ^2`$ between the data and a model spectrum based on a prescribed $`\tau (v)`$ can greatly underestimate the real column density errors. For example, let us assume that at a given point the covering fraction is exactly 0.5. In this case $`I=0.5+e^\tau `$, where $`I`$ is the residual intensity and $`\tau `$ is the optical depth. Assuming $`I=0.531`$, we obtain from the above relationship $`\tau =3.5`$. However, even a very small $`\mathrm{\Delta }I=0.03`$ error will cause the following error in optical depth, $`\tau =3.5_{0.8}^{+3.9}`$, and of course if $`\mathrm{\Delta }I=0.031`$ the formal error will be $`\tau =3.5_{0.8}^+\mathrm{}`$. An error of $`\mathrm{\Delta }I=.03`$ corresponds to a S/N=30, which is rarely (if ever) achieved for high resolution UV spectroscopy of AGN outflows. We therefore conclude that when the derived optical depth value is above 3 it is virtually impossible to set a meaningful positive error estimate. In such cases we should simply acknowledge that the derived optical depth value is a lower limit.

One can argue that a proper error analysis based on minimizing $`\chi ^2`$ should also yield the $`\tau =3.5_{0.8}^+\mathrm{}`$ result. However, this is only correct for one resolution element. Once we use a prescribed $`\tau (v)`$ to fit a full trough, we lose our ability to distinguish between highly saturated regions in the trough and less saturated ones. The formal $`\chi ^2`$ $`\tau `$ error will be obtained from averaging over the entire trough, whereas in some parts of the trough the actual error can be very large and dominate the derived column density error. This situation is accentuated when the shape of the trough is dominated by velocity dependent covering fraction and not by $`\tau (v)`$ variations.

For example, absorption component 4 in the Brotherton et al. (2002) modeling has a maximum optical depth of 7 in the blue doublet component (extracted from the parameters in their Table 2 and using their equations (1) and (2)), and therefore a maximum optical depth of 3.5 in the red doublet component. Based on the analysis presented above, it is possible that the real optical depth in this absorption component is far larger.

### 3.4 Different continuum emitting regions in the X-ray and UV

For O vi the relevant UV continuum is near 1030 Å, whereas that of the X-ray one is near 20 Å. It is plausible that the size and/or location of the emitting region at these wavelengths is quite different. It is possible that the outflow may cover different fractions of the continuum source, depending upon the wavelength. Such differences will complicate the comparison between UV and X-ray determined absorption column density of the same ion. Unfortunately, little is yet understood about the continuum emitting mechanism(s) in AGN. Thus, we cannot quantify the effects expected from the difference in the respective emitting regions.

However, in the context of the current investigation the important question is whether size/location differences in the emitting regions can explain the discrepancy between the UV and X-ray inferred O vi column densities. Two combined characteristics of the NGC 5548 O vi absorber leads us to conclude that such an outcome is not plausible. First, the velocity width and the overall depth of the O vi UV and X-ray absorption troughs are quite similar. Second, we are comparing an apparent O vi X-ray column density (derived from Eq. (1) using the apparent optical depth, i.e., assuming no saturation) with the inferred real UV O vi column density (which should account for saturation effects). It is easy to hide a large column density in a saturated trough (saturation in partially covering components or in unresolved components), but it is difficult to decrease the lower column density limit for a given trough’s shape. For example, if the covering fraction of the X-ray trough is much larger than the UV one, then the apparent X-ray column is larger than the apparent UV column. However, in that case we should observe a much deeper X-ray trough compared with the UV one, which is not seen in the NGC 5548 data. We conclude that it is unlikely that different UV and X-ray emitting region geometry will allow the apparent X-ray column density to be an order of magnitude larger than the real UV column density. Moreover, the striking kinematic similarity of the Ly$`\beta `$ UV trough and the O vii X-ray troughs in NGC 3783 (see Fig. 2 in Gabel et al. 2003) strongly suggests that a large overlap exists between the X-ray and UV continuum emitting regions along our outflow-intersecting line of sight.

We note though, that the argument can be turned around. A detailed comparison of X-ray and UV absorption troughs rising from the same ion may give us an important handle on the geometries of these different emitting regions. However, such a comparison necessitates higher quality X-ray data than those described here.

## 4 DISCUSSION

The O vi data in the NGC 5548 FUSE observation have a S/N$``$10 (Brotherton et al. 2002). Such limited S/N data do not allow a direct determination of the covering fraction and optical depth as a function of velocity. This forces the modeler to use some combination of physically weak assumptions, in descending order of “weakness” these are:

1. The apparent optical depths are equal to the real ones.
2. The optical depth is assigned a prescribed function (normally a Gaussian).
3. Each component of the flow is given a single covering fraction value, ignoring the possibility of velocity dependent covering.

Higher S/N observations are therefore needed in order avoid the practical necessity of using these assumptions.

We point out that a careful solution of the optical depth and covering fraction as a function of velocity (using doublet lines, e.g., Arav et al. 2002) can greatly increase the dynamical range for $`\tau `$ determination. Using the example from § 3.3, an $`I=0.531`$ gives $`\tau _{apparent}=0.6`$. Solving for $`\tau `$ using the doublet equations can yield robust estimates up to $`\tau =3`$ in the weak doublet component, or $`\tau =6`$ in the strong one, which increases our dynamical range for $`\tau `$ determination by an order of magnitude. Furthermore, even in cases of higher saturation, lower limits for $`\tau `$ and hence for $`N_{ion}`$ are useful in constraining the physical conditions of the absorber, provided we acknowledge that these are indeed lower limits (see the case of BALQSO PG 0946+301; Arav et al. 2001b).

Finally, in view of the caveats described in this paper, we advocate the following cautious recipe for determining column densities from absorption troughs of AGN outflows:

a) For a singlet line, with no other observed lines from the same ion, we can only measure the apparent optical depth as a function of velocity and use the resultant integrated column density as a lower limit. This is because it is conceptually impossible to decouple the effects of optical depth from those of covering factor for a single line, In particular, this is the case for Ly$`\alpha `$ in the HST/STIS band.

b) For a doublet line we can solve for both the optical depth and the covering fraction as a function of velocity. For this we need a sufficiently narrow absorption complex so that absorption features from the two doublet component do not blend. If blending occurs we have to resort to apparent optical depth estimates. Furthermore, for the unblended case, if the derived optical depth values are larger than $``$3 for the weak doublet component, we have to establish a realistic lower limit (based on the S/N and emission models uncertainties), which in practice will never be much higher than $`\tau `$3.

c) A line series from the same ion offer the best column density diagnostics, because of the large spread in oscillator strength values. Most important in this regard are the Lyman series. A modest S/N coverage of several of these lines allows for a simple and accurate determination of $`N_{H\mathrm{i}}`$. A similar situation occur in the X-ray band where many of the observed ions are H or He like and therefore show the corresponding line series.

## 5 SUMMARY

X-ray and UV spectra of NGC 5548 show O vi absorption troughs associated with the AGN outflow from this galaxy. A lower limit on the column density of the O vi X-ray trough is an order of magnitude larger than the column density found by Brotherton et al. (2002) using Gaussian modeling of the O vi UV troughs.

It is unlikely that the discrepancy in the inferred $`N_{O\mathrm{vi}}`$ reflect actual changes in the O vi absorption between the epochs of the X-ray and UV observations. Instead, we interpret the discrepancy as additional evidence that column densities inferred for UV troughs of Seyfert outflows are often severely underestimated. We identify the physical limitations of Gaussian modeling as the probable explanation of the $`N_{O\mathrm{vi}}`$ discrepancy. Specifically, Gaussian modeling cannot account for a velocity dependent covering fraction, and it is a poor representation for absorption associated with a dynamical outflow. Analysis techniques that use a single covering fraction value for each absorption component suffer from similar limitations.

## ACKNOWLEDGMENTS

This work is based on observations obtained with XMM-Newton, an ESA science mission with instruments and contributions directly funded by ESA Member States and the USA (NASA). Support for this work was provided by NASA through grant number HST-AR-09079 from the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., under NASA contract NAS5-26555 and by NASA LTSA grant 2001-029. NA expresses gratitude for the hospitality of the Astronomy department at the University of California Berkeley. The National Laboratory for Space Research at Utrecht is supported financially by NWO, the Netherlands Organization for Scientific Research. We also thank Yuri Levin for a thorough reading of the manuscript.


text:
# One-dimensional Dirac oscillator in a thermal bath

## Abstract

We analyze the one-dimensional Dirac oscillator in a thermal bath. We found that the heat capacity is two times greater than the heat capacity of the one-dimensional harmonic oscillator for higher temperatures.

The name Dirac oscillator was first introduced by Moshinky and Szczepaniak szczep  for a Dirac equation in which the momentum $`𝐩`$ is replaced by $`𝐩𝐩i`$m$`\omega \beta 𝐫`$ with $`𝐫`$ being the position vector, m the mass of the particle and $`\omega `$ the frequency of the oscillator. The name was given because the determination of the spectrum and eigenstate of a Dirac oscillator szczep  was obtained from that of an ordinary oscillator with a very strong spin-orbit coupling term. The concept, though not the name, had been discussed previously first by Ito et al. ito  and later by Cook cook  and Ui and Takeda Ui:zg  and possibly others. Besides intrinsic mathematical interest the study of the Dirac oscillator has invoked much attention because of its various physical applications. For instance, the last three chapter of a book of Moshinsky and Smirnov mosh  deal with the Dirac oscillator relativistic interactions between systems of one, two and three particles. Moreno and Zentella showed moreno  that it could be the object of an exact Foldy-Wouthuysen transformation. Benitez et al. Benitez:te  found the electromagnetic potential associated with the Dirac oscillator, and showed that this exactly soluble problem has a hidden supersymmetry, which is responsible for the special properties of its energy spectrum. They also calculated the related superpotential and discussed the implications of this supersymmetry on the stability of the Dirac sea.

Rozmej and Arvieu Rozmej:1999jv  have shown a very interesting analogy between the relativistic Dirac oscillator and the Jaynes-Cummings model. They showed that the strong spin-orbit coupling of the Dirac oscillator produces the entanglement of the spin with the orbital motion similar to what is observed in the model of quantum optics.

More recently Nogami and Toyama nogami  and Toyama et al toyama  have studied the behaviour of wave packets of the Dirac oscillator in the Dirac representation in (1+1) dimensions. The aim of these authors was to study wave packets which could possibly be coherent. This reduction of the dimension was brought as an attempt to get rid of spin effects and to concentrate on the relativistic effects.

As a matter of fact, the one-dimensional relativistic Dirac equation has been largely used to treat physical problems where relativistic effects could play an important role. Indeed, this treatment gives enlightening information about more realistic models in solid state physics solid state , particularly in semiconductor theories Renan:1999tr .

In spite of the great number of papers that has been recently published concerning the solution and properties of the Dirac oscillator, as far as we know no one has reported on its thermodynamics properties. In order to overcome this lack of information, in this letter we study the one-dimensional Dirac oscillator in a thermal bath.

Let us consider the energy spectrum for a one-dimensional Dirac oscillator adame

$$E_n=\sqrt{[2(n+1)]\mathrm{}\omega mc^2+m^2c^4},$$
(1)
where $`m`$ is the mass of the particle, $`\omega `$ is the angular frequency of the oscillator and $`n`$ is a positive integer number. As we can see, for $`n+1<<mc^2/2\mathrm{}\omega `$, the spectrum of the one-dimensional Dirac oscillator is approximated, up to $`E_0`$, to the spectrum of the one-dimensional harmonic oscillator.

The partition function of the Dirac oscillator at temperature $`T`$ is obtained through the Boltzmann factor,

$$Z=\underset{n=0}{\overset{\mathrm{}}{}}e^{(E_nE_0)\beta }=e^{\sqrt{b}\beta }\underset{n=0}{\overset{\mathrm{}}{}}e^{\sqrt{an+b}\beta },$$
(2)
were $`\beta =1/kT`$, $`k`$ the Boltzmann constant, $`a=2\mathrm{}\omega mc^2`$, $`b=m^2c^4+2\mathrm{}\omega mc^2`$, and we are using a frame in the heat bath. We have subtracted the ground state energy in order to make a comparison with the non-relativistic harmonic oscillator. All thermodynamics quantities for the Dirac oscillator are obtained through the partition function (2).

Let us first analyze the convergence of the partition function. The function $`f(x)=\mathrm{exp}(\sqrt{ax+b}\beta )`$ is a monotonically decreasing function and the corresponding integral

$$I(\beta )=_0^{\mathrm{}}e^{\sqrt{ax+b}\beta }𝑑x=\frac{2}{a\beta ^2}e^{\sqrt{b}\beta }(1+\beta \sqrt{b}),$$
(3)
is convergent. Thus from the theorems of convergent series, this implies that the partition function is convergent.

In order to evaluate the partition function, we can use the Euler-MacLaurin formula

$$\underset{n=0}{\overset{\mathrm{}}{}}f(n)=\frac{1}{2}f(0)+_0^{\mathrm{}}f(x)𝑑x\underset{p=1}{\overset{\mathrm{}}{}}\frac{1}{(2p)!}B_{2p}f^{(2p1)}(0),$$
(4)
where $`B_{2p}`$ are the Bernoulli numbers, $`B_2=1/6,B_4=1/30,\mathrm{}`$. They are defined through the series

$$\frac{t}{e^t1}=\underset{n=0}{\overset{\mathrm{}}{}}B_n\frac{t^n}{n!}.$$
(5)
Then the partition function can be written as

$$Z=\frac{1}{2}+\frac{2}{a\beta ^2}(1+\beta \sqrt{b})e^{\sqrt{b}\beta }\underset{p=1}{\overset{\mathrm{}}{}}\frac{B_{2p}}{(2p)!}f^{(2p1)}(0).$$
(6)
To compute the partition function, we need to calculate the sum in the above expression. For our case it can be done only by numerical methods. Up to $`p=2`$ this sum can be written as:

$`e^{\sqrt{b}\beta }{\displaystyle \underset{p=1}{\overset{2}{}}}{\displaystyle \frac{B_{2p}}{(2p)!}}f^{(2p1)}(0)=`$
$`\left({\displaystyle \frac{1}{24}}{\displaystyle \frac{\beta a}{\sqrt{b}}}{\displaystyle \frac{1}{1920}}{\displaystyle \frac{\beta a^3}{b^{5/2}}}{\displaystyle \frac{1}{1920}}{\displaystyle \frac{\beta ^2a^3}{b^2}}{\displaystyle \frac{1}{5760}}{\displaystyle \frac{\beta ^3a^3}{b^{3/2}}}\right).`$ (7)
Before computing numerically the partition function, let us analyze the case of higher temperatures. This correspond to $`\beta 1`$. All terms in the sum of (6) have a positive power in $`\beta `$, which are very small compared with the term $`2(1+\beta \sqrt{b})/a\beta ^2`$. Hence, for $`\beta 1`$

$$Z\frac{2}{a\beta ^2},$$
(8)
Now we can easily obtain the mean energy and the heat capacity of the Dirac oscillator for higher temperatures

$$U=\frac{\mathrm{ln}Z}{\beta }2kT,$$
(9)
$$C\frac{U}{T}2k.$$
(10)
These results show that, for higher temperatures, the mean energy and the heat capacity for the Dirac oscillator is two times the mean energy and the heat capacity of the non-relativistic one dimensional harmonic oscillator. This is due, only to relativistic effects, since in one dimension there is no spin coupling.

Now we briefly discuss our numerical results on the calculation of the partition function (6). First of all, we should mention that, using (6), the curves for the three thermal functions, namely, mean energy, heat capacity and free energy, are identical to those obtained for the nonrelativistic one-dimensional harmonic oscillator, in the limit of low temperatures ($`0`$ to $`400K`$).

We display a comparison between the Dirac oscillator and the harmonic oscillator only, for high temperatures (Fig. 1 and Fig. 2). It is seen that the free energy and the mean energy are greater for the Dirac oscillator than for the harmonic oscillator. From Fig. 3 it is also seen that the heat capacity for the Dirac oscillator is two times the heat capacity for the harmonic oscillator, result that was anticipated by the analytical calculations presented above.

It is worthwhile to mention that numerical calculations of the partition function for high temperatures ($`\beta 1`$), force us to handle several very small variables, which could lead to deceiving results. So, it is required a precise estimation of the involved physical quantities. Therefore, we choose the angular frequency of the oscillator as $`\omega =10^{20}Hz`$ in the region of high temperatures. We used adimensional quantities in the figures. The temperature ranges from $`10^8K`$ to $`10^{10}K`$ ($`0.01`$ to $`1.0`$ $`MeV`$).

As an extension of this work, we are currently studying a chain of Dirac oscillators in a thermal bath. Also, construction of higher dimensional Dirac oscillator in finite temperature, despite its technical difficulties, would be very interesting and could shed light on relativistic effects in statistical and solid state physics.

Some models used the nonrelativistic harmonic oscillator potential for describing confinement of quarks in mesons and baryons Ravndal:sr . More recently, some authors suggested moreno ; adame  that the Dirac oscillator could be a good candidate to be used as the confinement potential in heavy quark systems. On the other hand, recently ultrarelativistic heavy ion experiments have search for the quark gluon plasma, a novel phase of QCD in which quarks and gluons are deconfined Braun-Munzinger:2000mv . We expect that studies on quark-gluon plasma models and their thermal properties could be subsidized by our results.

ACKNOWLEDGMENTS

M. H. P., R. R. L. and C. A. S. A. were supported in part by Conselho Nacional de Desenvolvimento Científico e Tecnológico-CNPq.

REFERENCES