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# Some multidimensional integrable cases of nonholonomic rigid body dynamics 11footnote 1Journal Ref.: Regular & Chaotic Dynamics, 2003, 8 (1), 125-132.

## 1 The Suslov problem

The equations of nonholonomic systems are not Hamiltonian. They are Hamiltonian with respect to the almost-Poisson brackets . This is the reason that the integration theory of constrained mechanical systems is less developed then for unconstrained ones. However, in some solvable nonholonomic systems with an invariant measure, the phase space is foliated by invariant tori, placing these systems together with integrable Hamiltonian systems (see ).

In this paper we are interested in the integrability of multidimensional generalizations of the Suslov nonholonomic rigid body problem.

Consider an $`n`$–dimensional rigid body motion around the fixed point $`O=(0,0,\mathrm{},0)`$ in $`^n`$. The configuration space is the Lie group $`SO(n)`$. The matrix $`gSO(n)`$ maps the orthonormal frame $`E_1,\mathrm{},E_n`$ fixed in the space to the frame $`e_1=gE_1,\mathrm{},e_n=gE_n`$ fixed in the body, $`e_1=(1,0,\mathrm{},0)^t,\mathrm{},e_n=(0,\mathrm{},0,1)^t`$. For a motion $`g(t)SO(n)`$, the angular velocity and momentum (in body coordinates) are $`\mathrm{\Omega }(t)=g^1g(t)`$ and $`M=J(\mathrm{\Omega })`$ respectively. $`J:so(n)so(n)^{}`$ is the inertia tensor and has the form: $`M=J(\mathrm{\Omega })=I\mathrm{\Omega }+\mathrm{\Omega }I`$, where $`I`$ is symmetric $`n\times n`$ matrix called mass tensor (see ). Here we identified $`so(n)`$ and $`so(n)^{}`$ by the Killing scalar product.

Suppose that additional left–invariant constraints

$$a^i,\mathrm{\Omega }=0,i=1,\mathrm{},r$$
(1)
and potential force $`V(\mathrm{\Gamma })`$, where $`\mathrm{\Gamma }`$ is a unit vector fixed in the space, are given. Let $`D=\{\mathrm{\Omega }so(n),a^i,\mathrm{\Omega }=0,i=1,\mathrm{},r\}`$ and $`𝒟`$ be the corresponding left–invariant distribution of $`TSO(n)`$. The distribution $`𝒟`$ is nonintegrable (i.e., constraints are nonholonomic) if and only if $`D`$ is not a subalgebra of $`so(n)`$.

The equations of motion of a rigid body subjected to the constraints (1) and the potential force $`V(\mathrm{\Gamma })`$ are

$$\dot{M}=[M,\mathrm{\Omega }]+\frac{V}{\mathrm{\Gamma }}\mathrm{\Gamma }+\underset{i=1}{\overset{r}{}}\lambda _ia^i,\dot{\mathrm{\Gamma }}=\mathrm{\Omega }\mathrm{\Gamma },$$
(2)
where Lagrange multipliers $`\lambda _i`$ are determined from the constraints (1). To describe a motion $`g(t)`$ on $`SO(n)`$ one should solve an additional equation $`\mathrm{\Omega }=g^1g`$.

For $`n=3`$, $`r=1`$, the equations (1), (2) are equations of classical Suslov problem (see ). Therefore, the above system can be considered as a natural multidimensional generalization of the Suslov problem.

If $`D`$ is an eigenspace of $`J`$ then the equations (2) take the form

$$J(\dot{\mathrm{\Omega }})=\mathrm{pr}_D[J(\mathrm{\Omega }),\mathrm{\Omega }]+\mathrm{pr}_D(\frac{V}{\mathrm{\Gamma }}\mathrm{\Gamma }),\dot{\mathrm{\Gamma }}=\mathrm{\Omega }\mathrm{\Gamma },\mathrm{\Omega }D$$
(3)
Contrary to the general case (2), system (3) preserves the standard measure $`d\mathrm{\Omega }d\mathrm{\Gamma }`$ on $`=D\{\mathrm{\Omega }\}\times S^{n1}\{\mathrm{\Gamma }\},S^{n1}=\{\mathrm{\Gamma }^n,\mathrm{\Gamma },\mathrm{\Gamma }=1\}.`$ The energy $`E=\frac{1}{2}J\mathrm{\Omega },\mathrm{\Omega }+V(\mathrm{\Gamma })`$ is always an integral of (2), (3). If, besides the Hamiltonian function, we have $`dimD+n4`$ additional independent integrals, by the Jacobi theorem the solutions of (3) can be found by quadratures. Under the compactness conditions, the phase space $``$ is foliated by invariant tori with quasi-periodic dynamic (after appropriate change of time), similarly as in the Liouville theorem .

Therefore, it is natural to call the system completely integrable if it can be integrated by the Jacobi theorem; or more generally (see ): if the trajectories of the system belong to invariant tori with the dynamic of the form

$$\dot{\phi }_1=\frac{\omega _1}{\mathrm{\Phi }(\phi _1,\mathrm{},\phi _k)},\mathrm{},\dot{\phi }_k=\frac{\omega _k}{\mathrm{\Phi }(\phi _1,\mathrm{},\phi _k)}.$$
(4)
The flow of (4) in the new time $`d\tau =\mathrm{\Phi }^1dt`$ is quasi-periodic: $`d\phi _i/d\tau =\omega _i`$. Moreover, for almost all frequencies $`\omega _1,\mathrm{},\omega _k`$, by smooth change of variables $`\theta _i=\theta _i(\phi )`$ equations (4) can be reduced to the form $`\dot{\theta }_i=\mathrm{\Omega }_i=\omega _i/\mathrm{\Pi }`$, $`i=1,\mathrm{},k`$, where $`\mathrm{\Pi }`$ denotes the total measure of $`𝕋^k`$ (see ).

Note that the integrability of equations $`(\text{2})`$ without the constraints (for example the multidimensional Euler , Lagrange  and Klebsh  cases) implies non-commutative integrability of unconstrained system. The phase space $`T^{}SO(n)`$ is foliated on $`ddimso(n)`$ dimensional invariant isotropic tori with quasi-periodic dynamics (see ). In the nonholonomic case there is no Poisson structure. In order to precisely describe dynamic on the whole phase space $`𝒟`$ we have to solve kinematic equation $`g^1(t)\dot{g}(t)=\mathrm{\Omega }(t)`$. This problem, for the Fedorov–Kozlov integrable case, is studied by Zenkov and Bloch .

In the next section we give a review of the known integrable cases in three dimension. Then we present their natural higher dimensional generalizations. It appears that these systems are also completely integrable (section 3).

## 2 Three–dimensional integrable cases

For $`n=3`$ we can identify the Lie algebra $`(so(3),[,])`$ with $`(^3,\times )`$ via $`\mathrm{\Omega }=(\mathrm{\Omega }_1,\mathrm{\Omega }_2,\mathrm{\Omega }_3)=(\mathrm{\Omega }_{23},\mathrm{\Omega }_{13},\mathrm{\Omega }_{12})`$. Then the equations of the motion (2) become:

$$J\dot{\mathrm{\Omega }}=J\mathrm{\Omega }\times \mathrm{\Omega }+\mathrm{\Gamma }\times \frac{V}{\mathrm{\Gamma }}+\lambda a,\dot{\mathrm{\Gamma }}=\mathrm{\Gamma }\times \mathrm{\Omega },a,\mathrm{\Omega }=0,$$
(5)
Without loss of generality we can take $`a=(0,0,1)`$, i.e., the constraint is $`\mathrm{\Omega }_3=\mathrm{\Omega }_{12}=0`$. Then $`D=^2\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}`$. Geometrically, this means that only infinitesimal rotations in the planes $`e_1e_3`$ and $`e_2e_3`$ are allowed.

If $`V=0`$, then the equations (5) form a closed system in variables $`\mathrm{\Omega }_1,\mathrm{\Omega }_2`$. If $`a`$ is not an eigen vector of $`J`$ then there is an asymptotic line $`l`$ in $`D`$. The angular velocity is moving along the energy ellipse $`=\{E(\mathrm{\Omega })=h\}`$ asymptotically to intersections points $`w_{}`$ and $`w_+`$ of $``$ with $`l`$ (see Suslov ). That system, obviously, has not an invariant measure. When $`a`$ is an eigenvector of $`J`$, then system (5) preserves the measure $`d\mathrm{\Omega }d\mathrm{\Gamma }`$ and for $`V=0`$ the solutions are very simple $`\mathrm{\Omega }=const`$.

Thus, further we assume that $`a`$ is an eigenvector of the inertia tensor $`J`$. We can take $`J=\mathrm{diag}(J_1,J_2,J_3)`$. For the integrability, we need one additional integral, independent of the Hamiltonian function.

The known integrable cases are:

* The Lagrange case (see Kozlov ), when $`J_1=J_2`$ and $`V=B_3\mathrm{\Gamma }_3`$. The integral is $`F=J\mathrm{\Omega },\mathrm{\Gamma }`$.
* The Kharlamova case  where $`V(\mathrm{\Gamma })=B_1\mathrm{\Gamma }_1+B_2\mathrm{\Gamma }_2`$. Then $`F=J_1\mathrm{\Omega }_1B_1+J_2\mathrm{\Omega }_2B_2`$.
* The Klebsh–Tisserand case (Kozlov ) with $`V(\mathrm{\Gamma })=\frac{ϵ}{2}J\mathrm{\Gamma },\mathrm{\Gamma }`$ and $`F=\frac{1}{2}J\mathrm{\Omega },J\mathrm{\Omega }\frac{1}{2}A\mathrm{\Gamma },\mathrm{\Gamma }`$, where $`A=ϵJ^1detJ`$.

Let us note that in the first case the Lagrange multiplier is $`\lambda =0`$, and this is actually an invariant subsystem of the usual Lagrange rigid body system. Also, Ziglin proved that if $`V(\mathrm{\Gamma })=B_3\mathrm{\Gamma }_3`$ than system (5) has a complementary meromorphic first integral in the complexified phase space $`_{}=^2\{\mathrm{\Omega }_1,\mathrm{\Omega }_2\}\times S_{}^2\{\mathrm{\Gamma }\}`$ only in the Lagrange case $`J_1=J_2`$ .

The Kharlamova and Klebsh-Tisserand cases can be uniquely described and generalized in the following way (see ). If $`J_1J_2`$ then the equations (5) of the Suslov problem are integrable for the potentials: $`V(\mathrm{\Gamma })=v_1(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2^2+\mathrm{\Gamma }_3^2)+v_2(\mathrm{\Gamma }_2,\mathrm{\Gamma }_1^2+\mathrm{\Gamma }_3^2)`$ where $`v_1`$ and $`v_2`$ are arbitrary functions of two variables. The corresponding third integral is: $`F=\frac{1}{2}J\mathrm{\Omega },J\mathrm{\Omega }+J_2v_1(\mathrm{\Gamma }_1,\mathrm{\Gamma }_2^2+\mathrm{\Gamma }_3^2)+J_1v_2(\mathrm{\Gamma }_2,\mathrm{\Gamma }_1^2+\mathrm{\Gamma }_3^2).`$

In addition, there are integrable cases of the Suslov problem, which include the gyroscopic force with the momentum $`ϵ\mathrm{\Gamma }\times \mathrm{\Omega }`$ added to the right hand side of (5). Since the gyroscopic force is conservative, the Hamiltonian remains to be the first integral. As well as the usual Suslov problem, the system with the gyroscopic force is integrable for the potentials $`V=0`$, $`V(\mathrm{\Gamma })=B_1\mathrm{\Gamma }_1+B_2\mathrm{\Gamma }_2`$ and $`V(\mathrm{\Gamma })=\frac{1}{2}(a_1\gamma _1^2+a_2\gamma _2^2+a_3\gamma _3^2)`$ (see ).

## 3 Multidimensional integrable cases

Without potential force, the equations (1), (2) represent a closed system in variable $`\mathrm{\Omega }`$ (Euler–Poincare–Suslov equations). We can consider Euler–Poincare–Suslov equations on other Lie algebras as well. They are reduced equations of nonholonomic systems on Lie groups with left–invariant kinetic energy and left–invariant constraints (see ). Integrable examples can be found in .

Since in the three–dimensional case, only infinitesimal rotations in the planes $`e_1e_3`$ and $`e_2e_3`$ are allowed, Fedorov and Kozlov suggested that multidimensional analogue of Suslov’s conditions can be define in the following way: only infinitesimal rotations in the planes $`e_1e_n,\mathrm{},e_{n1}e_n`$, i.e., in the planes containing the vector $`e_n`$ are allowed . Therefore, we have the following constraints imposed to the components of the angular velocity matrix:

$$\mathrm{\Omega }_{ij}=0,i,jn1.$$
(6)
Then

$$=^{n1}\{\mathrm{\Omega }_{1n},\mathrm{\Omega }_{2n},\mathrm{},\mathrm{\Omega }_{n1,n}\}\times S^{n1}\{\mathrm{\Gamma }\}.$$
(7)
Really, for $`V=0`$, the system has the same behavior as three–dimensional system. If $`D`$ is an eigen space of the inertia operator $`J`$ then the solutions are simply $`\mathrm{\Omega }=const`$. In general, all trajectories lying on the energy ellipsoid $`=\frac{1}{2}J\mathrm{\Omega },\mathrm{\Omega }=h`$ are double asymptotic: they tend to the diametrically opposed points $`w_+`$ and $`w_{}`$ as $`t\pm \mathrm{}`$ (see ).

With the choice of constraints (6) we shall see that three–dimensional integrable cases with invariant measure have their multidimensional analogues.

Suppose that the mass tensor is diagonal $`I=\mathrm{diag}(I_1,I_2,\mathrm{},I_n)`$. Then $`D`$ is an invariant subspace of $`J`$. The orthogonal complement of $`D`$ in $`so(n)`$ is $`so(n1)`$. Since $`(so(n),so(n1))`$ is a symmetric pair ($`[so(n1),so(n1)]so(n1)`$, $`[so(n1),D]D`$, $`[D,D]so(n1)`$) the equations (3) take the following simple form:

$$J(\dot{\mathrm{\Omega }})=\mathrm{pr}_D(\frac{V}{\mathrm{\Gamma }}\mathrm{\Gamma }),\dot{\mathrm{\Gamma }}=\mathrm{\Omega }\mathrm{\Gamma },\mathrm{\Omega }D,$$

or coordinately:

$`(I_i+I_n)\dot{\mathrm{\Omega }}_{in}={\displaystyle \frac{V}{\mathrm{\Gamma }_i}}\mathrm{\Gamma }_n\mathrm{\Gamma }_i{\displaystyle \frac{V}{\mathrm{\Gamma }_n}},`$
$`\dot{\mathrm{\Gamma }}_i=\mathrm{\Gamma }_n\mathrm{\Omega }_{in},i=1,2,\mathrm{},n1,`$
$`\dot{\mathrm{\Gamma }}_n=\mathrm{\Gamma }_1\mathrm{\Omega }_{1n}+\mathrm{\Gamma }_2\mathrm{\Omega }_{2n}+\mathrm{}+\mathrm{\Gamma }_{n1}\mathrm{\Omega }_{n1,n}.`$ (8)
### 3.1 The Lagrange case

Let the mass tensor be of the form $`I=\mathrm{diag}(I_1,\mathrm{},I_1,I_n)`$, the rigid body be placed in a homogeneous gravitational field and the center of mass be on the axes of the dynamical symmetry $`e_n`$. Then the potential is of the form $`V(\mathrm{\Gamma })=B_n\mathrm{\Gamma }_n`$ and system (8) can be seen as an invariant subsystem of the unconstrained $`n`$–dimensional rigid body motion:

$`(I_i+I_j)\dot{\mathrm{\Omega }}_{ij}=(I_iI_j){\displaystyle \underset{k}{}}\mathrm{\Omega }_{ik}\mathrm{\Omega }_{kj}=0,i,jn1,`$
$`(I_1+I_n)\dot{\mathrm{\Omega }}_{in}=(I_1I_n){\displaystyle \underset{k}{}}\mathrm{\Omega }_{1k}\mathrm{\Omega }_{kn}B_n\mathrm{\Gamma }_n,`$
$`\dot{\mathrm{\Gamma }}_i={\displaystyle \underset{k}{}}\mathrm{\Omega }_{ik}\mathrm{\Gamma }_k,i=1,\mathrm{},n.`$ (9)
The Hamiltonian equations (9) represent the natural generalization of the three–dimensional Lagrange rigid body system. Their integrability is proved by Beljaev . However, this does not give us directly the integrability on the invariant subspace $`M_{ij}=\mathrm{\Omega }_{ij}=0`$, $`i,jn1`$. In order to prove the integrability of the Suslov problem we shall use the symplectic reduction, instead of verifying independence of Beljaev’s integrals on the invariant subspace.

Consider the Lagrange rigid body system on the whole phase space $`T^{}SO(n)`$. It is invariant with respect to the right Hamiltonian action of the group $`SO(n1)`$ (rotations which fix symmetry axes $`e_n`$). The corresponding moment map $`\mathrm{\Phi }:T^{}SO(n)so(n)^{}`$ is given by functions $`M_{ij}`$, $`i,jn1`$ considered as the left invariant functions on $`T^{}SO(n)`$. One can easily prove:

###### Lemma 3.1

The reduced system on $`\mathrm{\Phi }^1(0)/SO(n1)=T^{}(SO(n)/SO(n1))=T^{}S^{n1}`$ is the spherical pendulum.

The spherical pendulum is a completely integrable system. Therefore, the invariant submanifold $`\mathrm{\Phi }^1(0)T^{}SO(n)`$ of the Lagrange system, as well as the rest of the phase space, is almost everywhere foliated by invariant tori with quasi-periodic dynamics. (More about the relationships between integrability and reductions one can find in .) The Suslov problem on $`𝒟TSO(n)`$, after Legendre transformation, coincides with the Lagrange system on $`\mathrm{\Phi }^1(0)`$. Thus we obtain:

###### Theorem 3.1

The Suslov nonholonomic rigid body problem with constraints (6), the mass tensor of the form $`I=\mathrm{diag}(I_1,\mathrm{},I_1,I_n)`$ and the center of the mass on the axes of the dynamical symmetry $`e_n`$, in a presence of a homogeneous gravitational field (the multidimensional Lagrange case) is completely integrable.

Note that the unconstrained system on $`T^{}SO(n)`$ is invariant with respect to the left Hamiltonian action of the group $`SO(n1)`$ (rotations which fix $`\mathrm{\Gamma }`$), for any inertia tensor $`J`$ and potential $`V(\mathrm{\Gamma })`$. In this sense, the Lagrange system admits two different $`SO(n1)`$ reductions. For $`n=3`$ this corresponds to the fact that two Eulerian angles, the angle of precession $`\psi `$ and the angle of pure rotations $`\phi `$, are cyclic variables. Elimination of the variable $`\phi `$ with the zero value of the conjugate momentum $`p_\phi =0`$ gives the spherical pendulum on $`S^2`$ (for example, see ).

### 3.2 The Kharlamova case

The natural multidimensional generalization of the Kharlamova case can be defined as the constrained rigid body motion (8) in a gravitational field with the center of mass laying in the plane spanned by $`e_1,e_2,\mathrm{},e_{n1}`$. Then the potential force is of the form $`V=B_1\mathrm{\Gamma }_1+B_2\mathrm{\Gamma }_2+\mathrm{}+B_{n1}\mathrm{\Gamma }_{n1}`$ and equations (8) become:

$`(I_i+I_n)\dot{\mathrm{\Omega }}_{in}=B_i\mathrm{\Gamma }_n,`$
$`\dot{\mathrm{\Gamma }}_i=\mathrm{\Gamma }_n\mathrm{\Omega }_{in},i=1,2,\mathrm{},n1,`$
$`\dot{\mathrm{\Gamma }}_n=\mathrm{\Gamma }_1\mathrm{\Omega }_{1n}+\mathrm{\Gamma }_2\mathrm{\Omega }_{2n}+\mathrm{}+\mathrm{\Gamma }_{n1}\mathrm{\Omega }_{n1,n}.`$ (10)
Note that the singular points of (10) are given with $`\mathrm{\Gamma }_n=\dot{\mathrm{\Gamma }}_n=0`$. Also, equations (10) have integrals

$$F_{i,j}(\mathrm{\Omega })=\frac{I_i+I_n}{B_i}\mathrm{\Omega }_{in}\frac{I_j+I_n}{B_j}\mathrm{\Omega }_{jn},i,j=1,2,\mathrm{},n1.$$

Therefore, it is natural to introduce new coordinates $`\omega _1,\mathrm{},\omega _{n1},\gamma _1,\mathrm{},\gamma _n`$ in the following way:

$`\omega _1={\displaystyle \frac{I_1+I_n}{B_1}}\mathrm{\Omega }_{1n},`$
$`\omega _i={\displaystyle \frac{I_i+I_n}{B_i}}\mathrm{\Omega }_{in}{\displaystyle \frac{I_1+I_n}{B_1}}\mathrm{\Omega }_{1n},i=2,\mathrm{},n1,`$
$`\gamma _1={\displaystyle \frac{I_1+I_n}{B_1}}\mathrm{\Gamma }_1,`$
$`\gamma _i={\displaystyle \frac{I_i+I_n}{B_i}}\mathrm{\Gamma }_i+{\displaystyle \frac{I_1+I_n}{B_1}}\mathrm{\Gamma }_1,i=2,\mathrm{},n1,`$
$`\gamma _n=\mathrm{\Gamma }_n.`$ (11)
Then system (10) get the simple form

$`{\displaystyle \frac{d\omega _1}{dt}}=\gamma _n,`$
$`{\displaystyle \frac{d\omega _i}{dt}}=0,i=2,\mathrm{},n1,`$
$`{\displaystyle \frac{\gamma _i}{dt}}=\gamma _n\omega _i,i=1,2,\mathrm{},n1,`$
$`{\displaystyle \frac{\gamma _n}{dt}}=\left({\displaystyle \frac{B_1}{I_1+I_n}}\right)^2\gamma _1\omega _1{\displaystyle \underset{i=2}{\overset{n1}{}}}\left({\displaystyle \frac{B_i}{I_i+I_n}}\right)^2(\gamma _1+\gamma _i)(\omega _1+\omega _i).`$ (12)
Now, we can easily integrate the above system. We have $`\omega _2=\omega _2^0,\mathrm{},\omega _{n1}=\omega _{n1}^0.`$ By substitution of the first equation $`\dot{\omega }_1=\gamma _n`$ to the equations with $`\gamma _i`$–s we get the trajectory as a function of the variable $`\omega _1`$:

$`\gamma _1\gamma _1^0={\displaystyle \frac{1}{2}}((\omega _1)^2(\omega _1^0)^2),`$
$`\gamma _i\gamma _i^0=\omega _i^0(\omega _1\omega _1^0),i=2,\mathrm{},n1,`$ (13)
where $`\gamma _1^0,\mathrm{},\gamma _{n1}^0`$ are initial conditions.

Since $`\mathrm{\Gamma },\mathrm{\Gamma }=1`$, we have

$$\gamma _n^2=\mathrm{\Gamma }_n^2=1\underset{i=1}{\overset{n1}{}}\mathrm{\Gamma }_i^2=1\left(\frac{B_1}{I_1+I_n}\right)^2\gamma _1^2\underset{i=2}{\overset{n1}{}}\left(\frac{B_i}{I_i+I_n}\right)^2(\gamma _1+\gamma _i)^2=P_4(\omega _1),$$

where $`P_4`$ is a four degree polynomial depending on initial conditions as parameters. Finally, to get the trajectory as a function of the time $`t`$ one should solve the elliptic integral

$$tt_0=_{\omega _1^0}^{\omega _1}\frac{dw_1}{\pm \sqrt{P_4(\omega _1)}}.$$

The curve $`(\gamma _1(\omega _1),\mathrm{},\gamma _{n1}(\omega _1))`$ lies within the ellipsoid

$$\mathrm{\Sigma }=\left\{\left(\frac{B_1}{I_1+I_n}\right)^2\gamma _1^2+\underset{i=2}{\overset{n1}{}}\left(\frac{B_i}{I_i+I_n}\right)^2(\gamma _1+\gamma _i)^2=1\right\}.$$

On $`\mathrm{\Sigma }`$ we have $`\gamma _n=\mathrm{\Gamma }_n=0`$. Therefore the variable $`\omega _1`$ varies in the interval $`[\xi _1,\xi _2]`$ where $`\xi _1`$ and $`\xi _2`$ are adjacent roots of the polynomial $`P_4`$ between which it takes positive values and $`(\gamma _1(\xi _i),\mathrm{},\gamma _{n1}(\xi _i))\mathrm{\Sigma }`$. Thus, the trajectory is either asymptotic or periodic with period

$$T=2_{\xi _1}^{\xi _2}\frac{d\omega _1}{\sqrt{P_4(\omega _1)}}.$$

We can summarize the above considerations in the following theorem

###### Theorem 3.2

The phase space (7) of the multidimensional Kharlamova system is almost everywhere foliated by invariant circles. The periodic motions can be expressed as elliptic functions of time $`t`$.

### 3.3 The Klebsh–Tisserand case

By multidimensional Klebsh–Tisserand case we mean the nonholonomic rigid body motion (8) in the presence of the quadratic potential $`V(\mathrm{\Gamma })=\frac{1}{2}(B_1\mathrm{\Gamma }_1^2+\mathrm{}+B_n\mathrm{\Gamma }_n^2)`$:

$`(I_i+I_n)\dot{\mathrm{\Omega }}_{in}=(B_1B_n)\mathrm{\Gamma }_i\mathrm{\Gamma }_n,`$
$`\dot{\mathrm{\Gamma }}_i=\mathrm{\Gamma }_n\mathrm{\Omega }_{in},i=1,2,\mathrm{},n1,`$
$`\dot{\mathrm{\Gamma }}_n=\mathrm{\Gamma }_1\mathrm{\Omega }_{1n}+\mathrm{\Gamma }_2\mathrm{\Omega }_{2n}+\mathrm{}+\mathrm{\Gamma }_{n1}\mathrm{\Omega }_{n1,n}.`$ (14)
The system (14) is similar to the Fedorov–Kozlov integrable case . We have integrals

$$F_i(\mathrm{\Omega },\mathrm{\Gamma })=(B_iB_n)\mathrm{\Gamma }_i^2+(I_i+I_n)\mathrm{\Omega }_{in}^2i=1,2,\mathrm{},n1.$$

Note that the Hamiltonian function can be expressed in terms of $`F_i`$–s: $`E=\frac{1}{2}(nB_n+F_1+\mathrm{}+F_{n1})`$.

The phase space (7) of multidimensional Klebsh-Tisserand system is foliated by invariant varieties

$$𝒯_c=\{F_1=c_1,\mathrm{},F_{n1}=c_n\}.$$

Suppose that $`B_i>B_n`$, $`i=1,\mathrm{},n1`$. Then the functions $`F_i`$ are positive and $`S_{c_i}^1=\{F_i=c_i\}`$ are circles in the planes $`^2\{\mathrm{\Omega }_{in},\mathrm{\Gamma }_i\}`$, $`i=1,2,\mathrm{},n1`$. The invariant surface $`𝒯_c`$ is a two–covering of

$$S_{c_1}^1\times \mathrm{}\times S_{c_{n1}}^1\{\mathrm{\Gamma }_1^2+\mathrm{}+\mathrm{\Gamma }_{n1}^21\}$$

determined by the function $`\mathrm{\Gamma }_n=\pm \sqrt{1(\mathrm{\Gamma }_1^2+\mathrm{}+\mathrm{\Gamma }_{n1}^2)}`$. The branching points of the covering correspond to the zeros of $`\mathrm{\Gamma }_n`$. In particular, if constants $`c_i`$ satisfy inequality

$$\frac{c_1}{B_iB_n}+\mathrm{}+\frac{c_{n1}}{B_{n1}B_n}<1,$$
(15)
then $`\mathrm{\Gamma }_n0`$ and $`𝒯_c`$ is the union of two disjoint ($`n1`$)–dimensional tori (or less dimensional tori if some of $`c_i`$–s are equal to zero).

The motion on $`𝒯_c`$ can be solved in term of the new time $`d\tau =\mathrm{\Gamma }_ndt`$. Define the angular variables $`\phi _i`$ on circles $`S_{c_i}^1`$ by putting

$$\mathrm{\Omega }_{in}=\sqrt{\frac{c_i}{I_i+I_n}}\mathrm{sin}\phi _i,\mathrm{\Gamma }_i=\sqrt{\frac{c_i}{B_iB_n}}\mathrm{cos}\phi _i.$$

Using (14), we obtain the following quasi-periodic dynamics on $`𝒯_c`$:

$$\frac{\phi _i}{d\tau }=\omega _i=\sqrt{\frac{B_iB_n}{I_i+I_n}},i=1,2,\mathrm{},n1.$$

The frequencies $`\omega _i`$ depend only of $`I_i`$ and $`B_i`$. If the trajectories are periodic on one torus, they are periodic on the rest of the tori as well. Also, in the original time, system (14) takes the form (4) $`\dot{\phi }_i=\omega _i/\mathrm{\Gamma }_n^1`$.

###### Theorem 3.3

If $`c_1,\mathrm{},c_{n1}`$ satisfy the inequality (15) then the invariant variety $`𝒯_c`$ is the disjoint union of two invariant tori. The multidimensional Klebsh–Tisserand system (14) on $`𝒯_c`$ is completely integrable.

### Acknowledgments

I would like to thank M. Radnović for the kind support. The research was partially supported by the Serbian Ministry of Science and Technology, Project 1643 – Geometry and Topology of Manifolds and Integrable Dynamical Systems.


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# Topological properties of the Cosmic Microwave Background polarization map

## Abstract

Rapid progress has been made in observations of the temperature anisotropies of the Cosmic Microwave Background (CMB). These observations have enabled cosmologists to characterize the state of the universe at recombination, and observational efforts are now being directed towards obtaining the polarization map of the CMB. Here we draw an analogy between the CMB polarization map and nematic liquid crystals, pointing out that the two have similar defects. Making use of known results in the theory of defect formation we predict the statistical distribution of defects in the CMB polarization map and provide a novel tool for probing the CMB at large angular scales.

A polarization map of the CMB will consist of headless vectors laid on the surface of last scattering. One can view the polarization map as rod-shaped molecules laid out on a two-dimensional sphere ($`S^2`$), the sky. This picture immediately connects the CMB polarization map with nematics where the molecules are also rod-shaped. The topological properties of nematics have been studied for many decades Mer79 . The very topological properties that are relevant to nematics will also apply to the CMB polarization map, as first examined in Ref. DolDorNovNov99 .

At first sight, it might appear that the analogy is doomed because neighboring molecules in nematics interact, while CMB polarization vectors do not. However, the processes at last scattering that produce the polarization are described by interacting field theory and so neighboring polarization vectors, just as neighboring liquid crystal molecules, tend to align. The CMB polarization map is just like a snapshot of a thin film of nematic liquid crystal.

The mathematical description of the CMB polarization mapKos99  in a local patch of the sky is in terms of a $`2\times 2`$ symmetric traceless matrix which we denote by $`𝐏`$:

$$𝐏=\left(\begin{array}{cc}Q& U\\ U& Q\end{array}\right)$$
(1)
We define a normalized “order parameter” as: $`𝐌=\frac{1}{2}𝐏/\sqrt{Q^2+U^2}`$. The matrix $`𝐌`$ can also be written in terms of a unit vector $`\widehat{n}`$ as:

$$𝐌_{ij}=\widehat{n}_i\widehat{n}_j\frac{1}{2}\delta _{ij}$$
(2)
So the space of all matrices $`𝐌`$ is given by the space ($`S^1`$) of all two-dimensional unit vectors but with the $`Z_2`$ identification under $`\widehat{n}\widehat{n}`$. The normalized field $`𝐌`$ therefore defines a unit vector field with no arrows, also known as a “line field.”

Since $`𝐌`$ lies in $`S^1/Z_2`$, a space that has incontractable loops, the line field will have singularities, also called “defects.” Some of these singularities are shown in Fig. 1. The radial and tangential singularities commonly discussed in the CMB literature and shown in Fig. 2 can be constructed by combining two of the fundamental singularities.

In Euclidean space, each singularity has a degree $`W`$ which is given by:

$$W_\mathrm{\Gamma }=\frac{1}{2\pi }_\mathrm{\Gamma }𝑑x^i\mathrm{Tr}(𝐌ϵ_i𝐌)$$
(3)
where $`\mathrm{\Gamma }`$ is a contour around the location of the singularity, and

$$ϵ\left(\begin{array}{cc}0& 1\\ 1& 0\end{array}\right)$$
(4)
In terms of the functions $`Q`$ and $`U`$, we simply have

$$W_\mathrm{\Gamma }=\frac{1}{2\pi }_\mathrm{\Gamma }𝑑x^i_i\alpha $$
(5)
where $`\alpha =(1/2)\mathrm{tan}^1(U/Q)`$. The fundamental defects shown in Fig. 1 have degrees $`\pm 1/2`$, while the defects in Fig. 2 have degrees $`\pm 1`$. In general we can take $`\mathrm{\Gamma }`$ to be any closed contour on which the polarization is everywhere non-vanishing, and then Eq. (5) tells us the number of singularities within that contour. (We assume that all the singularities are isolated.)

Equation (3) can also be used on the surface of a sphere provided we first parallel transport $`𝐌`$ along a geodesic from points on $`\mathrm{\Gamma }`$ to a point within $`\mathrm{\Gamma }`$. In other words, to compute the degree, we have to parallel transport $`𝐌`$ to a common tangent space.

The Poincare-Hopf theorem ChoDeW89  tells us that the algebraic sum of degrees of all singularities, also called the “index” of a vector field, on a compact manifold is equal to the Euler number of the manifold. By explicit evaluation we can show that the Poincare-Hopf result also applies to a line field on $`S^2`$. Hence the CMB polarization map must necessarily have index 2 since the Euler number of $`S^2`$ is 2.

The distribution of defects produced at a phase transition has received much attention over the past few decades Kib76 . A key result is that if we consider the length, $`L`$, of $`\mathrm{\Gamma }`$ to be much larger than the correlation scale, $`\xi `$, of the matrix $`𝐌`$, then we can write VacVil84

$$W_{\mathrm{\Gamma },rms}W_\mathrm{\Gamma }^2^{1/2}=c\left(\frac{L}{\xi }\right)^\beta $$
(6)
where $`c`$ is a system-dependent constant but the exponent $`\beta `$ is universal with the value $`\beta =1/2`$. The reasoning leading to this scaling is very simple. $`W_\mathrm{\Gamma }`$ is the sum of random fluctuations in $`\alpha `$ and the number of elements in the sum goes like $`L/\xi `$. Therefore the root mean square value of $`W_\mathrm{\Gamma }`$ is proportional to $`\sqrt{L/\xi }`$.

Satellite observations of the CMB wmap ; cobe  show that the temperature anisotropy correlation function, $`C(\theta )`$, falls off very steeply with increasing angular separation, and becomes very small above an angular scale $`\theta =\chi `$ a few degrees. If this lack of large-scale temperature correlations is indeed real, and not some misunderstood systematic effect, it is natural to expect that one should find a corroborating lack of correlation for polarization on all angular scales larger than $`\chi `$. After all, thermal and polarization fluctuations are expected to arise from similar processes during recombination. Hence, on angular scales larger than $`\chi `$, the matrix $`𝐌`$ would be essentially uncorrelated. This tells us that the areal density of defects on the sky will be approximately one per $`\chi ^210`$ square degrees. Further, we expect $`\beta =1/2`$ on angular scales larger than $`\chi `$. If the data shows $`\beta 1/2`$ or exponential decay with $`L`$, then that would indicate the presence of large-scale correlations in the CMB polarization. This observation would be equivalent to defect formation with “bias” in condensed matter systems.

If the statistical distribution of the polarization is Gaussian, then all correlation functions of the polarization matrix $`𝐏`$ can be calculated in terms of the two point correlation function. However, correlators such as $`W_\mathrm{\Gamma }^2`$ that are central to our discussion, cannot be written in terms of two point correlators of $`𝐏`$, since they are given by correlators of the normalized matrix $`𝐌`$ which involves dividing $`𝐏`$ by $`\sqrt{Q^2+U^2}`$. Alternately, $`W_\mathrm{\Gamma }^2`$ depends on the two point correlation function of the angle $`\alpha =(1/2)\mathrm{tan}^1(U/Q)`$ and this is not given by the two point correlation functions of $`U`$ and $`Q`$. Hence the information provided by the distribution of defects is complementary to that provided by the two point correlators of the polarization matrix usually discussed.

There is a global constraint on the CMB polarization map, namely that “B-modes” are absent or highly suppressed KamKosSte97 ; SelZal97 ; KamKosSte97b ; SelZal97b  and the reader may worry that this would affect the distribution of defects. However, the B-mode relates the variation of the polarization intensity, $`\sqrt{Q^2+U^2}`$, to the distribution of the unit vector $`\widehat{n}`$. The topology, however, does not depend on the polarization intensity, only on the direction $`\widehat{n}`$. Hence for a fixed distribution of $`\widehat{n}`$, we would need to find $`\sqrt{Q^2+U^2}`$ such that the B-modes vanish. One can show that even pure E-modes exhibit the defects shown in Figs. 1 and 2 Bunetal03 .

As polarized CMB radiation propagates from the last scattering surface to us, intervening effects such as weak lensing and Faraday rotation can affect the polarization map. As long as these effects are continuous, they can only move the locations of the existing defects or, create or destroy defect-antidefect pairs. This is similar in spirit to the coarsening of defects observed in condensed matter systems, though the actual dynamics of polarization defects is likely to be very different.

In conclusion, we have drawn an analogy between the CMB polarization map and nematic liquid crystals, classified the defects in the polarization map, and made a prediction for the scaling of the index of the CMB polarization map. Ongoing observations of the CMB will be able to find the exponent $`\beta `$ and compare it to the theoretical prediction of $`1/2`$. Further statistical measures of the distribution of defects can provide additional novel probes of the CMB.

###### Acknowledgements.

We thank Craig Copi, Dragan Huterer, Arthur Kosowsky, Harsh Mathur, Levon Pogosian, John Ruhl, Rafael Sorkin, Glenn Starkman, and Phil Taylor for discussions. This work was supported by the U.S. Department of Energy and the CWRU Office of the Provost.


text:
# Semiclassical analysis of a nonlinear eigenvalue problem and non analytic hypoellipticity

## 1. Introduction

We are interested in a family of operators of the type

(1.1)
$$H(x,D_x,\lambda )=\mathrm{\Delta }+(P(x)\lambda )^2,$$

where $`xP(x)`$ is a polynomial of degree $`m`$. In the study of the failure of analytic hypoellipticity, one approach consists in showing that the nonlinear eigenvalue problem

(1.2)
$$H(x,D_x,\lambda )v=0$$

has at least one non trivial solution $`(\lambda ,v)C\times 𝒮(^n)`$ ($`v0`$). This has been used by many authors including Baouendi-Goulaouic, Helffer, Christ , Hanges-Himonas , Chanillo  and quite recently by Chanillo-Helffer-Laptev , where the reader can find a more extensive list of references. All the results lead to the formulation of a conjecture by Trèves  giving a necessary and sufficient condition of analytic hypoellipticity extending . It could be a rather natural conjecture that, when $`xP(x)`$ is an homogeneous elliptic polynomial on $`^n`$ of order $`m>1`$, (1.2) has at least one non-trivial solution. This result is proved () for $`n=1`$ and for $`m2n`$ when $`n=2,3`$ (). Our aim is to provide a semiclassical approach to this problem. Our analysis concerns actually a more general class of operators of the form

(1.3)
$$\underset{j=1}{\overset{p}{}}D_{x_j}^2+\left(P(x_1,\mathrm{},x_p)D_{x_{p+1}}D_{x_{p+2}}\right)^2+\left(Q(x_1,\mathrm{},x_p)D_{x_{p+1}}\right)^2.$$

When specialized to the case (1.2), we obtain the following

###### Theorem 1.1.

.
Let $`n`$ be even. Let $`P`$ be any real polynomial of degree $`m2`$ such that the homogeneous part of degree $`m`$ is elliptic. Then the nonlinear eigenvalue problem (1.2) has at least one non trivial solution.

Its now standard corollary to analytic hypoellipticity (see for example ) is :

###### Corollary 1.2.

.
Let $`P`$ be elliptic homogeneous of degree $`m2`$ on $`^p`$, with an even $`p>0`$, then the operator on $`^{p+2}`$ :

$$P(x,D_x):=\underset{j=1}{\overset{p}{}}D_{x_j}^2+\left(P(x_1,\mathrm{},x_p)D_{x_{p+1}}D_{x_{p+2}}\right)^2,$$

is not analytic hypoelliptic at $`0`$.

###### Remark 1.3.

As we shall discuss more deeply in Section 5, G. Métivier already showed in , that $`P(x,D_x)`$ is not hypoanalytic in any open set. This is indeed a sum of squares with an odd number of linearly independent vector fields. But the corollary presented here and its proof give a stronger information at $`0`$.

Other examples in any dimension related to (1.3) are also given in Section 4. These examples give an additional light to the general conjecture of .

Recently, Chanillo-Helffer-Laptev () used Lidskii’s Theorem to prove the existence of nonlinear eigenvalues of (1.2) permitting to recover some known results for $`n=1`$ and to give new examples in dimension $`n2`$.

The proof of Theorem 1.1 is based on the semiclassical analysis combined with the ideas of Chanillo-Helffer-Laptev (). We follow closely the reduction of Chanillo-Helffer-Laptev which is recalled in Section 2. In this approach via Lidskii’s theorem, the existence of a nonlinear eigenvalue of (1.2) is reduced to prove, for some $`k`$, the non vanishing of the trace of the k-th power of a linear operator $`𝒟`$:

$$\mathrm{Tr}𝒟^k0.$$

Our approach is to introduce artificially a semiclassical parameter $`h`$ in this operator (see Section 3). Then, the existence of a nonlinear eigenvalue of (1.2) is reduced to prove that for some $`k`$ and for some $`h`$,

$$\mathrm{Tr}𝒟(h)^k0,$$

where $`𝒟(h)`$ is an $`h`$-pseudo-differential operator. In other words, one can say that, while Chanillo-Helffer-Laptev try to prove that, for some $`k`$,

(1.4)
$$\underset{j}{}\lambda _j^k0,$$

where $`\lambda _j`$ are the nonlinear eigenvalues of (1.2) (which can not be real under our assumptions), we want to find here $`h`$ and $`k`$ such that

(1.5)
$$\underset{j}{}(h^{m/(m+1)}\lambda _j+1)^k0.$$

It is standard (see ) that $`\mathrm{Tr}𝒟(h)^k`$ has a complete semiclassical expansion when $`h0`$, which is theoretically computable. The question becomes now to find a non-zero term in these expansions which we call the semiclassical criterion. The computation of the leading term given in Section 4 gives the complete answer for $`n`$ even. Although we do not have a general result for all $`n`$ odd, we believe that the semiclassical criterion given below can be used to show that for each given value of $`n`$ odd, (1.2) has at least one non trivial solution. The introduction of a semiclassical parameter $`h`$ permits us to overcome the difficulty related to the non commutativity of operators encountered in ().

## 2. Chanillo-Helffer-Laptev’s approach

We rewrite $`H(x,D_x,\lambda )`$ in the form

(2.1)
$$H(x,D_x,\lambda )=L2\lambda M+\lambda ^2,$$

with

(2.2)
$$L=\mathrm{\Delta }+P(x)^2,M=P(x).$$

The operator $`L`$ is invertible and its inverse is a pseudo-differential operator (see ). It is also easy to give sufficient conditions for determining whether the operator

(2.3)
$$A:=L^1$$

belongs to a given Schatten class (see the appendix in ). Then the initial problem is reduced to the spectral analysis of

(2.4)
$$(I2\lambda B+\lambda ^2A)u=0.$$

with

(2.5)
$$B=A^{1/2}PA^{1/2}$$

Chanillo-Helffer-Laptev () used a rather standard approach to transform this nonlinear spectral problem to a linear ones and then apply Lidskii’s Theorem to prove the existence of a non trivial solution.

We first recall the reduction to the linear spectral problem. It is easy to see that it is enough to show that the operator $`𝒟`$ defined by

(2.6)
$$𝒟:=\left(\begin{array}{cc}2B& A^{1/2}\\ A^{1/2}& 0\end{array}\right)$$

has a non zero eigenvalue $`\mu `$. The first component of the corresponding eigenvector is an eigenvector of the problem (1.2) with $`\mu =\frac{1}{\lambda }`$.

If $`B`$ and $`A`$ are compact, $`𝒟`$ is compact but the main difficulty is that $`𝒟`$ is not selfadjoint. The Lidskii Theorem says that

###### Theorem 2.1.

.
Let $`𝒞`$ be a trace class operator. Let $`\lambda _j(𝒞)`$ denote the sequence of non zero eigenvalues of $`𝒞`$. Then

$$\underset{j}{}\lambda _j(𝒞)=\mathrm{Tr}𝒞.$$

Here the eigenvalues are repeated according to their algebraic multipicity.

By Lidskii’s Theorem, if $`\mathrm{Tr}𝒞0`$, $`𝒞`$ has at least one non zero eigenvalue.

###### Corollary 2.2.

(Chanillo-Helffer-Laptev ).
Let $`𝒟`$ be a compact operator. Assume that there exists $`k𝐍^{}`$ such that $`𝒟^k`$ is of trace class and $`\mathrm{Tr}𝒟^k0.`$ Then (1.2) has at least one non trivial solution.

Chanillo-Helffer-Laptev used this criterion in the case $`k=2,3,4`$ and obtained a family of interesting results in non analytic hypoellipticity. The non commutativity between $`A`$ and $`B`$ makes it hard to apply their method with higher rang criteria (although some results can be obtained in the same vein for $`k=6,8`$).

## 3. The semiclassical criterion

As explained in the introduction, our goal is to introduce a semiclassical parameter $`h`$ in the operators $`A`$ and $`B`$ and to apply the theory of $`h`$-pseudo-differential operators to give a complete semiclassical asymptotic expansion of the trace $`𝒟(h)^k`$. We are then led to find conditions under which the leading term is non zero.

Initially, the family of operators to study is of the form

(3.1)
$$H(x,D_x,\lambda )=\mathrm{\Delta }+(P(x)\lambda )^2,$$

where $`xP(x)`$ is a real polynomial on $`^n`$ of order $`m2`$. Write $`P`$ in the form

$$P=P_m+P_{m1}+\mathrm{}+P_0,$$

where $`P_j`$ is a homogeneous polynomial of degree $`j`$. We assume that $`P_m`$ is elliptic on $`^n`$:

(3.2)
$$P_m(x)0,x0.$$

To be definite, we assume $`P_m(x)>0,x0`$. The case $`P_m(x)<0,x0`$ is similar. This condition on $`P_m`$ requires $`m`$ to be even and thus excludes polynomials of odd degree when $`n=1`$.

To introduce the semiclassical parameter $`h`$, we make the dilation $`x\tau x`$ and set $`\lambda =(\lambda ^{}1)\tau ^m`$ (One takes $`\lambda =(\lambda ^{}+1)\tau ^m`$ if $`P_m`$ is negative). Let

(3.3)
$$H(x,hD;\lambda ^{},h)=h^2\mathrm{\Delta }+(P(x,h)\lambda ^{})^2,$$

where $`h=\frac{1}{\tau ^{m+1}}`$ and

(3.4)
$$P(x,h)=(P_m(x)+1)+h^{1/(m+1)}P_{m1}(x)+h^{2/(m+1)}P_{m2}(x)+\mathrm{}+h^{m/(m+1)}P_0.$$

Then, the initial problem

(3.5)
$$H(x,D;\lambda )v=0$$

has a non trivial solution $`(\lambda ,v)`$ if and only if

(3.6)
$$H(x,hD;\lambda ^{},h)u=0$$

has a non trivial solution $`(\lambda ^{},u)`$ for some $`h>0`$. Let us remark that

$$h^2\mathrm{\Delta }+(P_m+1)^21.$$

One can then prove that for $`h>0`$ small enough, $`h^2\mathrm{\Delta }+P(x,h)^2>1/2`$ and therefore is invertible.

More generally, let us consider the nonlinear eigenvalue problem

(3.7)
$$(h^2\mathrm{\Delta }+(Q(x,h))^2+(P(x,h)z)^2)u=0,$$

where $`P(x,h)`$ is of the above form and

(3.8)
$$Q(x,h)=Q_m(x)+h^{1/(m+1)}Q_{m1}(x)+h^{2/(m+1)}Q_{m2}(x)+\mathrm{}+h^{m/(m+1)}Q_0,$$

with $`Q_j`$ a homogenous polynomial of degree $`j`$. This kind of operators appears in the study of non analytic hypoellipticity for operators of the form (1.3). We assume that $`P_m`$ and $`Q_m`$ are real and there exists $`C>0`$ such that

(3.9)
$$C^1<x>^{2m}(P_m(x)+1)^2+Q_m(x)^2C<x>^{2m},x^n.$$

Here $`<x>=(1+|x|^2)^{1/2}`$. In the case $`Q_m=0`$, $`P_m`$ has to be an elliptic polynomial of degree $`m`$ so that (3.9) is satisfied. We can show by the argument used above that $`h^2\mathrm{\Delta }+P(x,h)^2+Q(x,h)^2`$ is invertible for $`h>0`$ small. We define the operators $`A=A(h)`$ and $`B=B(h)`$ in this new setting by :

(3.10)
$$A(h)=(h^2\mathrm{\Delta }+P(x,h)^2+Q(x,h)^2)^1,B(h)=A(h)^{1/2}P(x,h)A(h)^{1/2}.$$

For a temperate symbol, possibly $`h`$-dependent, $`K(x,\xi ,h)`$, we denote by $`K(x,hD,h)`$ the $`h`$-pseudo-differential operators defined by Weyl quantization

(3.11)
$$K(x,hD,h)u(x)=\frac{1}{(2\pi h)^n}e^{i(xy)\xi /h}K(\frac{x+y}{2},\xi ,h)u(y)𝑑y𝑑\xi ,$$

$`u𝒮(^n)`$. By the $`h`$-pseudo-differential calculus , we deduce easily that $`A(h)`$ and $`B(h)`$ are $`h`$-pseudo-differential operators with symbols, satisfying for any $`NIN^{}`$,

(3.12) $`a(x,\xi ;h)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{(m+1)N}{}}}h^{\frac{j}{(m+1)}}a_j(x,\xi )+h^{N+\frac{1}{(m+1)}}R_N(a,h)`$
(3.13) $`b(x,\xi ;h)`$ $`=`$ $`{\displaystyle \underset{j=0}{\overset{(m+1)N}{}}}h^{\frac{j}{(m+1)}}b_j(x,\xi )+h^{N+\frac{1}{(m+1)}}R_N(b,h).`$
Here the symbols

(3.14)
$$\begin{array}{cc}a_0\hfill & =(\xi ^2+(P_m+1)^2+Q_m(x)^2)^1,\hfill \\ b_0\hfill & =(\xi ^2+(P_m+1)^2+Q_m(x)^2)^1(P_m+1),\hfill \end{array}$$

are, by definition, the $`h`$-principal symbols of $`A(h)`$ and $`B(h)`$, respectively. If we denote by $`S_{\varphi ,\phi }^\rho `$ the class of symbols defined as in Robert , and if we introduce :

$$\begin{array}{cc}\rho _j\hfill & =(1+\xi ^2+x^{2m})^1(1+x^2)^{\frac{j}{2(m+1)}}(1+\xi ^2)^{\frac{j}{2(m+1)}},\hfill \\ \varphi \hfill & =(1+x^2)^{1/2},\hfill \\ \phi \hfill & =(1+\xi ^2)^{1/2},\hfill \end{array}$$

then

$$a_jS_{\varphi ,\phi }^{\rho _j},b_jS_{\varphi ,\phi }^{\rho _j<x>^m}.$$

Moreover the remainders $`R_N(a,h)`$ and $`R_N(b,h)`$ are respectively a bounded family of symbols in $`S_{\varphi ,\phi }^{\rho _{N(m+1)+1}}`$ and in $`S_{\varphi ,\phi }^{\rho _{Nm+1}<x>^m}`$.

The Chanillo-Helffer-Laptev’s $`k`$-criterion becomes now the following

###### Lemma 3.1.

.
Let $`𝒟(h)`$ be defined as in (2.6) with $`A`$, $`B`$ replaced by $`A(h)`$, $`B(h)`$. Let us assume there exists $`k`$ such that $`𝒟(h)^k`$ is an operator of trace class and that

(3.15)
$$\mathrm{Tr}𝒟(h)^k0$$

for some $`h>0`$. Then the nonlinear spectral problem (3.7) has at least one non trivial solution.

To apply this lemma, we prove the following

###### Theorem 3.2.

.
Let us assume the condition (3.9) for $`P`$ and $`Q`$ with $`m1`$. Let $`k>n(m+1)/m`$, $`n1`$. Then $`𝒟(h)^k`$ is an operator of trace class for all $`h>0`$ sufficiently small. We have, for any $`N`$, the following asymptotic expansion

(3.16)
$$\mathrm{Tr}𝒟(h)^k=(2\pi h)^n\{\underset{j=0}{\overset{N(m+1)}{}}h^{\frac{j}{(m+1)}}H_{j;n,k}+𝒪(h^{N+\frac{1}{(m+1)}})\},$$

when $`h0`$. Here $`H_{j;n,k}`$ is independent of $`h`$ and $`N`$, and can be computed from the symbol of $`𝒟(h)^k`$. In particular,

(3.17)
$$H_{0;n,k}=_{^{2n}}\mathrm{tr}\left(\sigma _k(x,\xi )\right)𝑑x𝑑\xi ,$$

$`\sigma _k`$ being the $`h`$-principal symbol of $`𝒟(h)^k`$.

Proof.
We note that $`A(h)^{1/2}`$ and $`B(h)`$ are $`h`$-pseudo-differential operators with symbol in $`S_{\varphi ,\phi }^{\rho _0^{1/2}}`$. Therefore $`𝒟(h)^k`$ is an $`h`$-pseudo-differential operator with matrix-valued symbol in the class $`S_{\varphi ,\phi }^{\rho _0^{\frac{k}{2}}}`$. Since $`m1`$, we have :

$$\rho _0^{\frac{k}{2}}L^1(^{2n})\text{ if }k>n(m+1)/m.$$

Consequently, $`𝒟(h)^k`$ is a trace class operator when $`k>n(m+1)/m`$. If we denote by $`\sigma _k(x,\xi ;h)`$ the total symbol of $`𝒟(h)^k`$, it has a complete semiclassical expansion similar to $`a(h)`$ beginning with $`\sigma _k`$, the $`h`$-principal symbol of $`𝒟(h)^k`$. The semiclassical expansion of the trace follows from the formula

$$\mathrm{Tr}𝒟(h)^k=(2\pi h)^n\mathrm{tr}\left(\sigma _k(x,\xi ;h)\right)𝑑x𝑑\xi .$$

Here $`\mathrm{tr}`$ denotes the trace of $`2\times 2`$ matrices.

A consequence of Lemma 3.1 and Theorem 3.2 is the following

###### Corollary 3.3.

(The semiclassical criterion).
Let (3.2) be satisfied. If there exist $`k`$ with $`k>(m+1)n/m`$ and $`j`$ such that $`H_{j;n,k}0`$, then the nonlinear eigenvalue problem (3.7) has at least one non trivial solution $`(z,u)`$ for each $`h>0`$ sufficiently small.

## 4. An application of the semiclassical criterion: the classical criterion

In this section, we apply the semiclassical criterion at the classical level, that is for $`j=0`$.

###### Proposition 4.1.

.
Assume that $`Q_m=0`$ and $`P_m0`$ is elliptic. Let $`k>(m+1)n/m`$. Then one has

(4.1) $`H_{0;n,k}`$ $`=`$ $`0,\text{ if }n\text{ is odd,}`$
(4.2) $`H_{0;n,k}`$ $`=`$ $`2(1)^{\mathrm{}}C_n{\displaystyle \frac{(n1)!}{(k1)(k2)\mathrm{}(kn)}}C(P_m),\text{ if }n=2\mathrm{}\text{.}`$
Here $`C_n`$ is the volume of $`S^{n1}`$ and

$$C(P_m)=_{_x^n}(P_m(x)+1)^k𝑑x>0.$$

In particular, we observe that :

(4.3)
$$H_{0;n,k}0\text{ when }k>(m+1)n/m$$

for all $`n`$ even. As a consequence, we get Theorem 1.1.

Proof of Proposition 4.1.
The condition (3.9) is satisfied. We compute

(4.4)
$$\mathrm{tr}\left(\sigma _k(x,\xi )\right)𝑑x𝑑\xi .$$

where $`\sigma _k`$ is the $`h`$-principal symbol of $`𝒟(h)^k`$. Although the symbolic calculus for matrix-valued $`h`$-pseudo-differential operators is complicated, the $`h`$-principal symbol of $`𝒟(h)^k`$ can be easily computed. Since the $`h`$-principal symbol of $`𝒟(h)`$ is

$$\left(\begin{array}{cc}2b_0& a_0^{1/2}\\ a_0^{1/2}& 0\end{array}\right)$$

the $`h`$-principal symbol of $`𝒟(h)^k`$ is

$$\sigma _k=\left(\begin{array}{cc}2b_0& a_0^{1/2}\\ a_0^{1/2}& 0\end{array}\right)^k.$$

Therefore,

(4.5)
$$\text{ tr }\sigma _k=(b_0+\sqrt{b_0^2a_0})^k+(b_0\sqrt{b_0^2a_0})^k.$$

We recall that

$$a_0=(\xi ^2+(P_m+1)^2)^1,b_0=(\xi ^2+(P_m+1)^2)^1(P_m+1).$$

By the change of variables $`\xi (P_m+1)\eta `$, we obtain that

(4.6)
$$\mathrm{tr}(\sigma _k)(x,\xi )𝑑x𝑑\xi =C(P_m)_{_\eta ^n}(1+\eta ^2)^k((1+i|\eta |)^k+(1i|\eta |)^k)𝑑\eta .$$

This shows that

(4.7)
$$H_{0;n,k}=2C_nC(P_m)\mathrm{}_0^{\mathrm{}}(1+r^2)^k(1+ir)^kr^{n1}𝑑r,k>(m+1)n/m,$$

$`C_n`$ being the volume of $`S^{n1}`$.

Let

$$L(n,k)=_0^{\mathrm{}}(1+r^2)^k(1+ir)^kr^{n1}𝑑r=_0^{\mathrm{}}(1ir)^kr^{n1}𝑑r.$$

An integration by parts gives :

$`L(n,k)`$ $`=`$ $`i{\displaystyle \frac{n1}{k1}}L(n1,k1)`$
$`=`$ $`i^{n2}{\displaystyle \frac{(n1)(n2)\mathrm{}2}{(k1)(k2)\mathrm{}(kn+2)}}L(2,kn+2).`$
Since

$$L(2,j)=\frac{1}{(j1)(j2)},$$

we have

(4.9)
$$H_{0;n,k}=2C_n\mathrm{}\{i^n\frac{(n1)!}{(k1)(k2)\mathrm{}(kn+1)(kn)}\}C(P_m).$$

This proves (4.1) and (4.2).

Using the classical criterion, we can construct a family of examples when $`n`$ is odd.

###### Corollary 4.2.

.
Let $`n=n_1+n_2`$ with $`n_1`$ even and $`n_21`$. Let $`P`$, $`R`$ be real elliptic homogeneous polynomials111The important property is only that : $`\mathrm{\Delta }_{x^{\prime \prime }}+R(x^{\prime \prime })^2`$ on $`^{n_2}`$ has an eigenvalue $`\lambda _0`$. of degree $`m`$ on $`^{n_1}`$ and $`^{n_2}`$, respectively. For $`x^n`$, set $`x=(x^{},x^{\prime \prime })^{n_1}\times ^{n_2}`$. Then, the operator :

$$\mathrm{\Delta }_x+\left(R(x^{\prime \prime })D_{y_1}\right)^2+\left(P(x^{})D_{y_1}D_{y_2}\right)^2,$$

is not analytic hypoelliptic at $`0`$ in $`^{n+2}`$.

Proof.
It suffices to show that the nonlinear eigenvalue problem

$$\left(\underset{j=1}{\overset{n}{}}D_{x_j}^2+\left(R(x^{\prime \prime })\right)^2+\left(P(x^{})z\right)^2\right)u=0$$

has a non trivial solution $`(z,u)`$. We can look for $`v(x^{})`$ satisfying, for some $`z`$,

$$\left(\mathrm{\Delta }_x^{}+\lambda _0+\left(P(x^{})z^{}\right)^2\right)v=0$$

where $`\lambda _0>0`$ is an eigenvalue of $`\mathrm{\Delta }_{x^{\prime \prime }}+R(x^{\prime \prime })^2`$. The corresponding semiclassical operator is $`h^2\mathrm{\Delta }_x^{}+\lambda _0h^{2m/(m+1)}+\left(P(x^{})z\right)^2`$ on $`^{n_1}`$. We can apply Proposition 4.1 for $`Q=Q_0=\lambda _0`$ to conclude.

###### Remark 4.3.

.
One can play with $`\lambda _0`$. We can indeed consider a sequence $`\lambda _n`$ of eigenvalues of $`\mathrm{\Delta }_{x^{\prime \prime }}+R(x^{\prime \prime })^2`$ tending to $`+\mathrm{}`$ and to choose a corresponding sequence $`h_n`$ of $`h`$’s tending to $`0`$, such that $`\lambda _nh_n^{\frac{2m}{m+1}}=1`$. This could permit to relax the condition of ellipticity of $`P`$, particularly if we add a $`Q`$ like in (3.7).

Let us now consider the more general case $`Q_m0`$. For $`(P_m,Q_m)`$ satisfying (3.9), we can check that the classical criterion still works in some cases. Let us remark that

(4.10)
$$\text{ tr }\sigma _k=(b_0+\sqrt{b_0^2a_0})^k+(b_0\sqrt{b_0^2a_0})^k$$

with now

$$a_0=(\xi ^2+(P_m+1)^2+Q_m^2)^1,b_0=(\xi ^2+(P_m+1)^2+Q_m^2)^1(P_m+1).$$

Let

(4.11)
$$T=((P_m+1)^2+Q_m^2)^{1/2},\tau _1=(P_m+1)/T,\tau _2=|Q_m|/T.$$

By the change of variables $`\xi T\eta `$, we obtain that

(4.12)
$$\mathrm{tr}(\sigma _k)(x,\xi )𝑑\xi =2T^{nk}\mathrm{}_{_\eta ^n}(1+\eta ^2)^k(\tau _1+i(\tau _2^2+\eta ^2)^{1/2})^k𝑑\eta .$$

#### The case $`n=1`$

.
This case can be also thoroughly analyzed by the method of . When $`n=1`$, we obtain here for $`k=3`$,

(4.13)
$$H_{0;1,3}=\frac{3\pi }{2}__x\tau _1\tau _2^2T^2𝑑x.$$

It is clear that if $`P_m0`$ (this requires $`m`$ to be even) and $`Q_m0`$, then, $`\tau _10`$ and $`H_{0;1,3}<0`$. If $`P_m`$ changes sign, $`H_{0;1,3}`$ may be vanishing.

#### The case $`n2`$

.
When $`n2`$, $`k>(m+1)n/m`$,

(4.14)
$$H_{0;n,k}=2C_n\mathrm{}_{_x^n}T(x)^{nk}_0^{\mathrm{}}(1+r^2)^k(\tau _1+i(\tau _2^2+r^2)^{1/2})^kr^{n1}𝑑r.$$

Set

(4.15)
$$C_{n,k}=_0^{\mathrm{}}(1+r^2)^k(\tau _1+i(\tau _2^2+r^2)^{1/2})^kr^{n1}𝑑r.$$

By a change of variable, we obtain

$$C_{n,k}=_{\tau _2}^{\mathrm{}}(\tau _1it)^k(t^2\tau _2^2)^{(n2)/2}t𝑑t.$$

#### The subcase $`n=2`$

.
In the case $`n=2`$, it is easy to see that

(4.16)
$$C_{2,k}=\frac{(\tau _1+i\tau _2)^{k2}}{k2}\frac{\tau _1(\tau _1+i\tau _2)^{k1}}{k1}.$$

Since $`\tau _1^2+\tau _2^2=1`$, we can see that

(4.17) $`\mathrm{}C_{2,3}`$ $`=`$ $`\tau _1(1+2\tau _2^2)/20,\text{ if }\tau _10`$
(4.18) $`\mathrm{}C_{2,4}`$ $`=`$ $`{\displaystyle \frac{1}{2}}+2\tau _1^2{\displaystyle \frac{4}{3}}\tau _1^4.`$
This shows $`H_{0;2,3}>0`$ for $`n=2`$ when $`P_m0`$. The classical criterion can be used in this case for $`m>3`$ when $`n=2`$.

To study the case $`n=2`$ and $`m=2`$, we need information on the sign of $`H_{0;2,4}`$ which depends on the relation between $`P_2`$ and $`Q_2`$. An elementary computation gives :

###### Lemma 4.4.

.

1. If $`\tau _1(x_1,x_2)^2(3\sqrt{3})/4`$ for all $`(x_1,x_2)`$, then $`H_{0;2,4}>0`$.
2. If $`\tau _1(x_1,x_2)^2(3\sqrt{3})/4`$ for all $`(x_1,x_2)`$, then $`H_{0;2,4}<0`$.

#### The subcase $`n=3`$

.
When $`n=3`$, $`k=5`$, we have

(4.19)
$$H_{0;3,5}=\frac{5\pi }{8}_{_x^3}\tau _2^4\tau _1T^2𝑑x>0,$$

so long as $`Q_m`$ is not identically zero and $`P_m0`$.

#### Some applications

.
We obtain the following :

###### Proposition 4.5.

.
(a). Let $`n=1`$. The operator

$$P_m(x,D):=D_x^2+(D_scx^mD_t)^2+(x^mD_t)^2$$

is not analytic hypoelliptic at $`0^3`$ for all $`m2`$ even and $`c`$.

(b). Let $`n=2`$. Let $`P`$, $`Q`$ be real homogenous polynomials of degree $`m2`$ satisfying (3.9) and $`P0`$. When $`m=2`$, we assume in addition that one of the above conditions on $`\tau _1`$ is satisfied. Then the operator

(4.20)
$$D_x^2+D_y^2+(D_sP(x,y)D_t)^2+(Q(x,y)D_t)^2,(x,y)^3,$$

is not analytic hypoelliptic at $`0`$ in $`^4`$.

(c). Let $`n=3`$. Let $`P`$, $`Q`$ be real homogenous polynomials on $`^3`$ of degree $`m2`$ satisfying (3.9). Assume that $`Q0,P0`$. Then the operator

(4.21)
$$\mathrm{\Delta }_x+(D_sP(x)D_t)^2+(Q(x)D_t)^2$$

is not analytic hypoelliptic at $`0`$ in $`^5`$.

We believe that the condition on $`\tau _1`$ in (b) is technical. As examples of $`(P,Q)`$ satisfying the conditions of (b) of Proposition 4.5, we can take $`P(x,y)=(x^2+y^2)^{\mathrm{}}`$, $`Q(x,y)=(xy)^{\mathrm{}}`$, $`\mathrm{}1`$, because in the case $`\mathrm{}=1`$, one can check that $`\tau _1(x,y)^24/5>(3\sqrt{3})/4`$. We shall come back at the analysis of this example corresponding to $`\mathrm{}=1`$ in the next section.

###### Remark 4.6.

.
Let us briefly compare the results obtained here and those of . Apparently, in the case $`Q=0`$, the “classical” criterion does not produce any result for $`n1`$ odd. But the semiclassical criteria may still work by considering a coefficient of higher order of the expansion of the trace. As indicated by J. Sjöstrand, a similar approach is used by L. Nedelec in  for getting a lower bound for the number of resonances of an $`h`$-pseudodifferential system. The same condition on the dimension appears. The “quantum” criterion given in  works for $`n=1,2,3`$ but with a stronger condition on $`m`$ when $`n>1`$. Moreover it has been observed in Remark 4.4 in  that the condition of ellipticity of $`P`$ can be replaced by a weaker condition. The last point is that the homogeneity of $`P`$ plays an important role in the dilation argument of , while in the semiclassical approach, the lower order parts can be included. This appears to be useful in the dimension reduction.

## 5. Comparison with Métivier’s results

In this section, we would like to analyze the links between our results and the previous results obtained by G. Métivier .

### 5.1. A first family of examples

Let us first consider the operator on $`^{n+2}`$ :

$$H(X,D_X)=\mathrm{\Delta }+(P(x)D_{x_{n+1}}D_{x_{n+2}})^2,$$

where $`X=(x,x_{n+1},x_{n+2})`$, $`P`$ is an homogeneous positive elliptic polynomial of degree $`m2`$ on $`^n`$.

If we take the “microlocal spirit” we observe that $`H`$ is an operator with double characteristics, whose principal symbol is the function :

$$(T^{}^{n+2}0)(X,\mathrm{\Xi })|\xi |^2+(P(x)\xi _{n+1}\xi _{n+2})^2.$$

The symbol vanishes exactly at order $`2`$ on the submanifold

$$\mathrm{\Sigma }=\{(X,\mathrm{\Xi })|\xi =0,P(x)\xi _{n+1}\xi _{n+2}=0,\xi _{n+1}0\}.$$

This submanifold is of codimension $`n+1`$. Let us now analyze the “symplecticity” of $`\mathrm{\Sigma }`$. We recall that $`\mathrm{\Sigma }`$ is said to be symplectic if the restriction to $`\mathrm{\Sigma }`$ of the canonical $`2`$-form is non degenerate. An easy way for verifying the symplecticity is to consider the $`(n+1)\times (n+1)`$ matrix $`u_i,u_j`$ where $`u_i(X,\mathrm{\Xi })=\xi _i`$ for $`i=1,\mathrm{},n`$ and $`u_{n+1}(X,\mathrm{\Xi })=P(x)\xi _{n+1}\xi _{n+2}`$ and to show that it is not degenerate. An immediate computation shows that its rank at a given point is $`2`$ if $`P0`$ and $`0`$ if $`P=0`$.

When $`P`$ is elliptic and homogeneous, we get that the rank is constant outside $`0`$ and equal to $`2`$. There are two cases :

1. When $`n=1`$, we get that $`\mathrm{\Sigma }`$ is symplectic except at the points of $`\mathrm{\Sigma }`$ such that $`x=0`$. The result of Trèves-Tartakoff-Métivier-Sjöstrand  shows that the operator is microlocally analytic hypoelliptic outside of $`\mathrm{\Sigma }`$ (ellipticity) and in the neighborhood of the points of $`\mathrm{\Sigma }`$ such that $`x0`$. In this case, the operator is not analytic hypoelliptic at any point $`(0,x_{n+1},x_{n+2})`$.
2. When $`n>1`$, Métivier’s result222Note that the operator is hypoelliptic with loss of one derivative in $`^{n+2}\{x=0\}`$.  gives that the operator $`H`$ is not analytic hypoelliptic in any open set in $`^{n+2}`$. What we show here is the more precise result that $`P`$ is not analytic hypoelliptic at any point $`(0,x_{n+1},x_{n+2})`$ which is a finer property. See the introduction of  for the discussion between the definitions of analytic hypoellipticity and germ-hypoanalyticity (analytic hypoellipticity in a neighborhood of a point and analytic hypoellipticity at a point).

### 5.2. A new class of non analytic hypoelliptic operators

We now show that maybe more interesting examples can be treated when considering the more general class :

$$H(X,D_X)=\mathrm{\Delta }+(P(x)D_{x_{n+1}}D_{x_{n+2}})^2+Q(x)^2D_{x_{n+1}}^2,$$

where $`P`$ and $`Q`$ are homogeneous polynomials of degree $`m>1`$ with $`P0`$ and $`P^2+Q^2`$ elliptic. When restricting $`H`$ to $`x_{n+1}`$ independent distributions, we get an analytic hypoelliptic operator on $`^{n+1}`$:

$$\mathrm{\Delta }_x+(P(x)^2+Q(x)^2)D_{x_{n+2}}^2,$$

by applying a theorem of Grushin ().

We have seen in Section 4, that the “classical” criterion can give a result under some additional condition. We will concentrate our analysis to the specific case when :

$$n=2,P(x)=x_1^2+x_2^2,Q(x)=\alpha x_1x_2,$$

with $`\alpha >0`$.

Let us do the same microlocal analysis as in the previous subsection. The characteristic set $`\mathrm{\Sigma }`$ is now defined as a union of two regular submanifolds of dimension $`4`$ in $`^80`$ :

$$\mathrm{\Sigma }=\mathrm{\Sigma }_1\mathrm{\Sigma }_2,$$

with

$$\mathrm{\Sigma }_j=\{\xi _1=0,\xi _2=0,\xi _4=(x_1^2+x_2^2)\xi _3,x_j=0,\xi _30\}.$$

Moreover $`\mathrm{\Sigma }_j`$ is symplectic outside $`\mathrm{\Sigma }_1\mathrm{\Sigma }_2`$ and not symplectic at $`\mathrm{\Sigma }_1\mathrm{\Sigma }_2`$.

Outside $`\mathrm{\Sigma }_1\mathrm{\Sigma }_2`$, observing that the symbol of $`H`$ vanishes exactly at order $`2`$ on $`\mathrm{\Sigma }`$ there, we get again from  that $`H`$ is microlocally analytic hypoelliptic. Métivier’s criterion of non-analytic-hypoellipticity can not be applied at the points $`(0,0,x_3,x_4)`$ (the operator is indeed not hypoelliptic with loss of one derivative) and it is interesting to see what is obtained through our approach.

###### Proposition 5.1.

.
For $`\alpha >0`$ small enough, the operator $`H(X,D_X)`$ is not analytic hypoelliptic at any point $`(0,x_3,x_4)`$.

Proof.
We just apply the criterion of the previous section (Proposition 4.5, second case) and the discussion following the statement.

When $`\alpha `$ is small enough (at least $`0<\alpha 1`$), we observe that the quantity $`(P_m+1)^2/\left((P_m+1)^2+Q_m^2\right)`$ is sufficiently near $`1`$. In particular, the second condition on $`\tau _1`$ in Lemma 4.4 is satisfied.

###### Remark 5.2.

.
As explained to us by M. Christ, it is possible to prove, for rather large classes of models depending analytically on an additional parameter $`\alpha `$, that some associated Fredholm determinant is analytic in $`\alpha `$. One can conclude that, if the operator is not analytic hypoelliptic for some value of $`\alpha `$, then it is not analytic hypoelliptic for generic values of $`\alpha `$. We refer to , Proposition 5.2 for the argument. In the particular case of the above Proposition 5.1, we can present the argument in the following way. Let $`H_{0,2,4}(\alpha )`$ denote $`H_{0,2,4}`$ defined as in Theorem 3.2 with $`P=x_1^2+x_2^2`$ and $`Q=\alpha x_1x_2`$. Using (4.12) and (4.16), one can check that $`H_{0,2,4}(\alpha )`$ is real analytic in $`\alpha >0`$. The proof of Proposition 5.1 shows that $`H_{0,2,4}(\alpha )0`$ for $`\alpha >0`$ small. Thus, $`H_{0;2,4}(\alpha )0`$ for all $`\alpha >0`$, except a discrete set in $`_+`$ and therefore Proposition 5.1 remains true in this case. This argument can be used for more general cases. It is indeed easy to check in many situations the analyticity of $`H_{0,n,k}`$ with respect to the parameter.


text:
# Current-voltage characteristics of quasi-one-dimensional superconductors: An S-curve in the constant voltage regime

## Abstract

Applying a constant voltage to superconducting nanowires we find that its IV-characteristic exhibits an unusual S-behavior. This behavior is the direct consequence of the dynamics of the superconducting condensate and of the existence of two different critical currents: $`j_{c2}`$ at which the pure superconducting state becomes unstable and $`j_{c1}<j_{c2}`$ at which the phase slip state is realized in the system.

The majority of the experiments on the resistive state in quasi-one dimensional systems were performed in the constant current regime and at temperatures close to $`T_c`$. It is extremely difficult to apply voltage to a superconductor because the current density induced by the applied electric field inevitably reaches the critical value and destroys superconductivity in the sample. The decrease of the superconducting current by the appearance of phase slip centers Skocpol ; Kramer1 ; Kramer2 ; Kramer3 ; Watts-Tobin  is not effective in this case because of the large heating of the sample at low temperatures. At temperatures close to $`T_c`$ the heating can be suppressed due to the low value of the critical currents but in this case the applied electric field does not penetrate deep into the sample because of the existence of regions near the N-S boundary where the drop of the applied voltage occurs Yu ; Ryazanov .

This situation drastically changes with the appearance of nano-technology and the ability to create long (to allow the appearance of phase slip centers) superconducting wires with a small cross section (to decrease the effect of heating). In this Letter we present results on the behavior of such nanowires in the constant voltage regime. We found that the I-V characterestic in this case has a remarkable S-shape. Our theoretical analysis based on the time-dependent Ginzburg-Landau equations (TDGL) shows that such a behavior is a direct consequence of the dynamics of the superconducting condensate and we predict new unusual features which still need additional experimental study.

The superconducting nanowires were prepared by electrodeposition into nanopores of homemade track-etched polycarbonate membranesLegras . For the lead nanowires, a 22 $`\mu `$m thick membrane (with pore diameter $``$ 40 nm and pore density $``$ 4$``$109 cm-2) and an aqueous solution of 40.4 g/l Pb(BF4)2, 33.6 g/l HBF4 and 15 g/l H3BO3 were used Michotte , while in the case of the tin nanowires, a 50 $`\mu `$m thick membrane (with pore diameter $``$ 55 nm and pore density $``$ 2$``$109 cm-2) and an electrolyte of 41.8 g/l Sn(BF4)2 in water solution were applied. Constant potential of -0.5 V versus an Ag/AgCl reference electrode was used in a three-electrode configuration in order to reduce the Pb2+ or Sn2+ ions into the nanopores. As shown in Fig. 1, the nanowires are cylindrical and the diameter is uniform along their length.

In order to perform electrical transport measurements on a single nanowire (still kept inside the membrane), a self-contacting technique was developed. In addition of the thick gold cathode ($``$ 1 $`\mu `$m) from which the nanowires start to grow up inside the membrane, another thin gold layer (in the range of 50 to 200 nm) was deposited on the other side exposed to the electrolyte prior to the electrodeposition. This contacting layer covers the surface but in contrast to the cathode, it is not thick enough to plug the pores. As the nanowires do not grow at the same speed, the first emerging nanowire interrupts the growth of the others by favoring the formation of a film of the deposited material on this contacting layer.

All the transport measurements are done in an electromagnetically shielded environment and special precautions were taken to reduce the noise. The voltage across the sample was measured by a Keithley voltmeter. From the numerous measurements performed both on Pb and on Sn nanowires, it appears that almost no depression of Tc was observed. In addition, a superconducting transition was found to occur when the residual resistivity ratio (RRR) is larger than about 2 while for smaller RRR values an insulating behavior is observed Michotte . Our nanowires are moderately disordered and the RRR of the samples was larger than 10. Structural characterization of the Pb nanowires using TEM experiments confirm also their polycrystalline character Dubois .

In Fig. 2, the DC current-voltage (I-V) characteristics of the Pb and Sn nanowires are reported, applying either a DC current or a DC voltage.

All the curves presented in Fig. 2 were measured in both directions and no hysteresis was observed. These measurements were also performed with reversed polarity without any changes compared to the reported results. In the usual current driven experiment we observe for the Pb nanowire (Fig. 2(a)) the successive appearance of two PSCs at low temperature. Applying the phenomenological SBT model Skocpol , the size of the first PSC can be estimated to be about 18 $`\mu `$m, which is twice the quasiparticles diffusion length. In spite that successive PCSs tend to avoid those already in place, the second PSC created in this Pb nanowire is thus forced to interpenetrate the first one, which explains why the second jump in resistance is smaller. The current driven experiment on the 50 $`\mu `$m long Sn nanowire (Fig. 2(b)) also shows the formation of a PSC at low temperature but for this particular material, the PSC extension is larger (around 40 $`\mu `$m) and the formation of a second PSC almost coincides with the transition to the normal state. Interesting new features were observed when measuring the reverse way, i.e. applying constant voltage and measuring the current. Here, the current flowing into the nanowire was determined by measuring the voltage across a 1 $`\mathrm{\Omega }`$ resistance added in series and the voltage developed across the sample was measured separately. In this voltage driven experiment, we found an interesting S-shape behavior which occurs at low temperature in the formation region of these PSCs, both for the Pb and the Sn nanowires.

To understand the experimental results we used the generalized TDGL equation Kramer1  which describes the dynamics of the superconducting condensate $`\psi =|\psi |e^{i\theta }`$

$$\frac{u}{\sqrt{1+\gamma ^2|\psi |^2}}\left(\frac{}{t}+i\phi +\frac{\gamma ^2}{2}\frac{|\psi |^2}{t}\right)\psi =\frac{^2\psi }{s^2}+\psi (1|\psi |^2).$$
(1)
The physical quantities are presented in dimensionless units (see e.g. Ref. Kramer2 ) and we neglected the vector potential, because in our case the self-induced magnetic field is negligible in size. Eq. (1) is supplemented with an equation for the electrostatic potential $`\phi `$ which is obtained from the condition for the conservation of the total current in the wire, i.e. $`\mathrm{div}𝐣=0`$. The parameter $`\gamma `$ depends on the material and we took $`u5.79`$ Kramer1 .

Although Eq. (1) is valid only close to $`T_c`$ (see Kramer1 ) we expect to find clues of the dynamics of the system even far from $`T_c`$. This expectation is motivated by the fact that results found from the solution of the stationary Ginzburg-Landau equations for mesoscopic superconductors Deo  agree with experiment even far from $`T_c`$.

If we apply an external current to the wire then the system exhibits hysteretic behavior Kramer1 ; Kramer3 . If one starts from the superconducting state and increases the current the superconducting state switches to the resistive superconducting or normal state at the upper critical current $`j_{c2}`$. With decreasing current it is possible to keep the sample in this state even for currents up to $`j_{c1}<j_{c2}`$ (which we call the lower critical current). The superconducting resistive state is realized as a periodic oscillation of the order parameter in time at one point of the superconductor Kramer1 ; Kramer2 . When the order parameter reaches zero in this point, a phase slip of $`2\pi `$ occurs. Such a state is now called a phase slip state and the corresponding point a phase slip center (PSC). The meaning of the current $`j_{c1}`$ is that for $`j<j_{c1}`$ the phase slip solution cannot be realized (in the absence of fluctuations) and thus the current $`j_{c1}`$ is the critical current at which PSC’s start to appear. Thus there is a limiting cycle (see Ref. Ivlev ) for Eq. (1) giving rise to this process.

The existence of such two critical currents is, in the $`V=const`$ regime, responsible for a complicated I-V characteristic (see Fig. 3). The I(V) curve has an overall S-like behavior which at low voltages exhibits an oscillatory behavior superimposed on it. We found that these unusual properties are typical for superconducting wires with $`\gamma 1`$. To understand the physics leading to this unusual behavior we consider the dynamics of $`\psi `$ in the wire for different voltages.

When we apply a voltage $`V`$ to the wire of length $`L`$ a constant electrical field appears in the sample $`E=j_E=V/L`$. As a result the superconducting condensate will be accelerated and the momentum $`p=\theta `$ will increase in time. For small voltages and electric fields the absolute value of the order parameter and current density are approximately described by quasi-equilibrium expressions with $`|\psi |(t)=1p(t)^2`$ and $`j_s(t)=p(t)(1p(t)^2)`$ ($`p(t)j_Et`$) while $`p<p_c`$ (at $`L/\xi 1`$ we have $`p_c1/\sqrt{3}`$). For $`p>p_c`$ the above spatial independent solution becomes ’unstable’ (like for the case of a ring in the presence of an applied magnetic field as discussed in Refs. Tarlie ; Vodolazov ). The order parameter vanishes in the center of the wire and a jump of the phase by $`2\pi `$ occurs in this point. It decreases $`p`$ by $`\mathrm{\Delta }p=2\pi /L`$ (and hence the superconducting current density) in the entire wire. A second phase slip will appear in the same point and the phase slip process repeats itself as long as $`j>j_{c1}`$. When the current density becomes less than $`j_{c1}`$ the order parameter will increase till its quasi-equilibrium value $`|\psi |(t)=1p(t)^2`$ with $`p(t)<p_c`$. Than the applied electric field will again accelerate the superconducting condensate and the above process is repeated (see Fig. 4). We define the time which is necessary to decrease the order parameter from $`|\psi |=\sqrt{2/3}`$ till zero and again back to the value $`|\psi |=\sqrt{2/3}`$ as the transition time $`T_{tr}`$ and call this entire process the transition process.

It is easy to show that at small voltages (in the limit $`V0`$) the number of phase slips which occurs during the transition process in an ideal wire is given by the following expression $`\mathrm{N}_{\mathrm{min}}=\mathrm{Nint}\left((\mathrm{p}_\mathrm{c}\mathrm{p}_{\mathrm{c1}})(\mathrm{L}/2\pi +1)\right)`$, where $`p_{c1}`$ is the smallest real root of equation $`j_{c1}=p_{c1}(1p_{c1}^2)`$ and $`\mathrm{Nint}(\mathrm{x})`$ returns the nearest integer value. The number of phase slips occurring during the transition process increases with increasing wire length.

Because of the presence of an external electric field, $`p`$ also increases during the transition process. Roughly speaking, if the voltage in the wire is so large that $`j_ET_{tr}>2\pi /L`$ the number of phase slips during the transition will increase by one (see Fig. 4). It leads to an instant decrease of the average current. The reason is that the maximum current during current oscillations $`j_{max}`$ for a small change of $`V`$ does not change but the minimum value of the current $`j_{min}`$ decreases by about $`2\pi /L`$, and hence $`<I>`$ should decrease. This effect is most pronounced in short wires where $`N_{min}`$ is small and the term $`2\pi /L`$ is large. If $`N`$ does not change with increasing $`V`$ then $`<I>`$ increases because both minimum and maximum currents grow.

When the period $`T=2\pi N/V`$ becomes of order $`2T_{tr}`$ the dynamics of the system changes considerably (at voltage $`V_1`$ in Fig. 3). Starting from this point the average current $`<I>`$ monotonically decreases with increasing applied voltage. For these values of the voltages the system is unable to return to the ”quasi-equilibrium” state (it is a state such that if we switch off the voltage instantly the system remains in this state for an infinitely long time in the absence of fluctuations). At any moment of time the order parameter is more suppressed in the point where the oscillations occur. At these voltages $`j_{max}`$ decreases with increasing $`V`$ and in principle it should reach the current $`j_{c1}`$ (and $`j_{min}`$ should reach $`j_{c1}`$ from the bottom). But this is only possible in an infinitely long wire. In a finite length wire every phase slip event leads to a decrease of the momentum $`p`$ and the current $`j`$ by a finite value. As a result the lower critical current in the $`V=const`$ regime is equal to $`j_{c1}^V=j_{c1}+dj`$ (with $`dj1/L`$) at which the ’usual’, like in the $`I=const`$ regime, periodic phase slip processes with period $`T=2\pi /<V>`$ starts. This corresponds to the voltage $`V_2`$ in Fig. 3.

In long wires a complex behavior is found in the current decreasing region due to the existence of several phase slip events during the transition process. At the initial parts of this region an irregular behavior is typical. In Fig. 5 we show (dotted curve) an example of such a situation (at $`V=0.08`$ \- point 4 in Fig. 3). Not a single, but several periods $`T=2\pi N/V`$ are present. This regime is replaced at higher voltages (for example at $`V=0.2`$ \- point 5 in Fig. 3) by a regime with a single period and the situation becomes similar to the case at low voltages.

We expect that the S shape of the I-V characteristic in the $`V=const`$ regime will be conserved even for temperatures far from $`T_c`$, because the only condition for its occurrence is the existence of two different critical currents $`j_{c1}<j_{c2}`$. While $`V<V_1`$ the current in the wire will be an oscillating function with maximal amplitude $`j_{c2}`$ and minimal amplitude $`j_{c1}`$. It means that the time averaged current should be larger than $`j_{c1}`$. At $`V>V_1`$ we expect that $`j_{max}`$ decreases due to the same mechanisms as was discussed above and hence the average current should decrease. But the maximal current cannot decrease below $`j_{c1}`$ and for $`V>V_2`$ both $`j_{min}`$ and $`j_{max}`$ have to increase with increasing voltage.

In our theoretical consideration we neglected the effect of thermo- and/or quantum-activated phase-slip centers. We believe that this is the reason why in the experiment no hysteresis was observed in the $`I=const`$ regime. The larger these fluctuations the larger they will decrease the measured $`j_{c2}^{exp}`$ and increase the measured $`j_{c1}^{exp}`$ (see also the discussion of this question in Ref. Kramer4 ) and in principle they may suppress the hysteresis. But in the $`V=const`$ regime the situation is quite different. If the time needed for the appearance of a fluctuating phase slip center is much larger than the period of oscillation of the current then we can neglect their effect.

In conclusion, we presented the first experiment on the I-V characteristic in the voltage constant regime for narrow superconducting Pb and Sn wires. An unexpected S-shape response is found which is absent in a current driven measurements. A theoretical analysis based on the generalized TDGL equations gives a qualitative explanation of this phenomena which is driven by a highly nonlinear time response of the superconductor to electric fields leading to the periodic creation of PS centers. Furthermore, we predict an oscillatory I-V characteristic at low voltages which was not found in the experiment presumedly due to the large length of the wire and hence the relatively large value of $`N_{min}`$.

This work was supported by IUAP(P5/1/1), GOA, ESF, FWO-Vl and the ”Communauté Française de Belgique” through the program ”Actions de Recherches Concertées”. We thank the POLY lab. at UCL for the polymer membranes. S. Mi. is a Research Fellow of the FNRS.


text:
# Origin of the Superconductivity in 𝛽'-(BEDT-TTF)2ICl2 under High Pressure and 𝛽-(BEDT-TTF)2X at Ambient Pressure

## Abstract

We present a theoretical study on the superconductivity in $`\beta ^{}`$-(BEDT-TTF)2ICl2 at $`T_\mathrm{c}=`$14.2K under high hydrostatic pressure found by Taniguchi et al., which is the highest record among organic superconductors. Its electronic structure is well expressed by the anisotropic triangular lattice Hubbard model at half filling. In the present work, we study this effective model by using the fluctuation-exchange (FLEX) approximation. In the obtained phase diagram, the superconductivity with $`d_{x^2\text{-}y^2}`$ like symmetry is realized next to the antiferromagnetic (AF) insulating phase, as a result of the 1D-2D dimensional crossover driven by the pressure. The obtained maximum $`T_\mathrm{c}`$ is 16$``$18K. In addition, the superconductivity in $`\beta `$-(BEDT-TTF)2X is also understood in the same framework.

PACS numbers: 74.70.Kn, 74.70.-b

Up to now, many kinds of organic superconductors composed of BEDT-TTF (abbreviated as ET) molecules have been discovered and studied intensively. Especially, $`\kappa `$-(ET)2X compounds attract much attention because it shows the maximum superconducting (SC) transition temperature ($`T_\mathrm{c}12`$K) among them . Recently, the temperature-pressure phase diagram of $`\kappa `$-(ET)2X salt could be produced by the fluctuation-exchange (FLEX) approximation . It is a kind of the self-consistent spin-fluctuation theory . According to the FLEX approximation, $`d`$-wave superconductivity is expected which is mediated by the strong antiferromagnetic (AF) fluctuations due to the Coulomb interaction.

Quite recently, Taniguchi et al. found the superconductivity at $`T_\mathrm{c}=14.2`$K in $`\beta ^{}`$-(ET)2ICl2 under high hydrostatic pressure (P$`\stackrel{>}{}`$8.2GPa), which established the new record of $`T_\mathrm{c}`$ among the organic superconductors . At ambient pressure, $`\beta ^{}`$-(ET)2ICl2 shows a semiconducting conductivity below the room temperature, and indicates a magnetic transition at $`T_\mathrm{N}`$=22K. As the pressure is applied, the resistivity decreases gradually, and the metallic behavior ($`d\rho /dT>0`$) is observed above $`T_{\mathrm{MIT}}`$ under 6.5GPa. Note that $`T_{\mathrm{MIT}}`$, which is a crossover temperature, is expected to be higher than $`T_\mathrm{N}`$. At 8.2GPa, insulating phase (or $`T_{\mathrm{MIT}}`$) disappears and the SC transition occurs at $`T_\mathrm{c}=14.2`$K at the same time.

According to the band calculation, the Fermi surface (FS) of $`\beta ^{}`$-(ET)2X (X=ICl2, BrICl, AuCl2) compound is quasi one-dimensional (Q1D) . It is contrastive that the FS of $`\beta `$-(ET)2X is almost round, which shows the superconducting transition at $`T_\mathrm{c}=2.7`$K for X=IBr2, at $`T_\mathrm{c}=3.8`$K for X=AuI2 and at $`T_\mathrm{c}=1.5`$K (or 8K) for X=I3. The difference of the crystal structure between $`\beta `$ and $`\beta ^{}`$ compounds arises from the fact that the anion size in $`\beta ^{}`$ compounds is too small to retain the $`\beta `$-type structure ; in fact, the bond lengths for I-Cl and Au-Cl are 2.52Å and 2.27Å respectively, which are smaller than the I-Br bond length in $`\beta `$-(ET)2IBr2, 2.58Å .

Until now, the deformation of the structure of $`\beta ^{}`$-(ET)2ICl2 at $`P=8.2`$GPa has not been determined experimentally. In the present study, however, we naturally assume that its structure (gradually) approaches to that of $`\beta `$-(ET)2X as the cell volume decreases under hydrostatic high pressure, because the anion size becomes larger compared with the cell length. Note that the shape of the unit cells for $`\beta `$\- and that for $`\beta ^{}`$-structures are almost same.

In the present work, we study the origin of the superconductivity in $`\beta ^{}`$-(ET)2ICl2 on the assumption that the crystal structure approaches to $`\beta `$-type under hydrostatic pressure. We use the FLEX approximation by noticing the fact that the SC phase appears by destroying the AF phase by pressure. Our theory can reproduce a reasonable value of $`T_\mathrm{c}`$ in $`\beta ^{}`$-(ET)2ICl2 under high pressure, as well as a low $`T_\mathrm{c}`$ in $`\beta `$-(ET)2X at ambient pressure.

The schematic structures of $`\beta ^{}`$-(ET)2X and $`\beta `$-(ET)2X are shown in Fig. 1. Each ellipse represents the ET molecule, and each highest occupied molecular orbital (HOMO) of the ET molecule possesses 1.5 electrons on average. The hopping parameters obtained by the band calculation are given in Table I: The hopping parameters for $`\beta `$’-(ET)2AuCl2 and those for $`\beta `$’-(ET)2ICl2 are similar. Whereas, they are very different from the parameters for $`\beta `$-(ET)2AuCl2 although their crystal structures differ slightly. This comes from the fact that the overlap integral between two ET molecules is very sensitive to the angle between them ($`\varphi `$): It takes the local maximum (positive) values for $`\varphi =0`$o,60o and the local minimum (negative) values for $`\varphi =30`$o,90o .

As a good approximation, we take account of only the anti-bonding orbit of each pair of the ET molecules connected by $`t(p1)`$, which is more than 2.5 times larger than other hopping parameters . Then, the original systems is mapped onto the “dimer model” at half filling, which is the anisotropic triangular lattice Hubbard model with four parameters ($`t_0`$, $`t_1`$, $`t_2`$,$`U_{\mathrm{eff}}`$) as shown in Fig. 1. Note that the $`x`$-$`y`$ coordinate in the present model does not correspond to that for the original systems. The corresponding hopping parameters for $`\beta ^{}`$-(ET)2ICl2 and for $`\beta `$-(ET)2I3 are given in Table II.

|  | $`\beta `$’-(ET)2ICl2 | $`\beta `$’-(ET)2AuCl2 | $`\beta `$-(ET)2I3 |
| --- | --- | --- | --- |
| $`t(p1)`$ | $`0.272`$ (eV) | $`0.264`$ (eV) | $`0.245`$ (eV) |
| $`t(p2)`$ | $`0.016`$ | $`0.020`$ | $`0.084`$ |
| $`t(q1)`$ | $`0.100`$ | $`0.100`$ | $`0.127`$ |
| $`t(q2)`$ | $`0.066`$ | $`0.065`$ | $`0.068`$ |
| $`t(c)`$ | $`0.016`$ | $`0.023`$ | $`0.050`$ |

Table I: Hopping integrals for each systems. They are given by the well known empirical relation $`t=10S`$ (eV), where $`S`$ is the overlap integral given by the band calculation .

|  | $`x=0`$ \[$`\beta `$’-(ET)2ICl2\] | $`x=1`$ \[$`\beta `$-(ET)2I3\] |
| --- | --- | --- |
| $`t_0`$ | $`0.050`$ (eV) | $`0.064`$ (eV) |
| $`t_1`$ | $`0.017`$ | $`0.084`$ |
| $`t_2`$ | $`0.008`$ | $`0.042`$ |

Table II: Hopping integrals in the dimer model for $`x=0`$ and $`x=1`$, respectively.

Although the original on-site Coulomb interaction on a ET molecule is $`1`$eV, the effective Coulomb interaction on a dimer, $`U_{\mathrm{eff}}`$, is limited to $`2|t(p1)|0.5`$eV .

Hereafter, we study the dimer model given in Fig. 1 at half filling. To analyze the pressure effect on the electronic states, we put the hopping parameters of the dimer model ($`t_0`$, $`t_1`$, $`t_2`$) as follows:

$`t_0=00.5+0.014x(\mathrm{eV}),`$ (1)
$`t_1=00.17+0.067x(\mathrm{eV}),`$ (2)
$`t_2=0.008+0.05x(\mathrm{eV}),`$ (3)
where $`0x1`$ is a free parameter. $`x=0`$ ($`x=1`$) corresponds to $`\beta ^{}`$-(ET)2ICl2 ($`\beta `$-(ET)2I3); see table II. Non-linear terms with respect to $`x`$ are dropped in eq.(2). The $`x`$ dependence of the FS for $`U_{\mathrm{eff}}=0`$ is shown in Fig. 2. We see that the FS becomes round around $`x0.7`$.

Here, we calculate the self-energy for the dimer model for $`U_{\mathrm{eff}}=0.40.5`$eV by using the FLEX approximation . Hereafter, we write $`U_{\mathrm{eff}}`$ as $`U`$ for simplicity. To obtain the magnetic transition temperature $`T_\mathrm{N}`$, we calculate the Stoner factor without vertex corrections (VC’s), $`\alpha _\mathrm{S}`$, given by

$`\alpha _\mathrm{S}=\underset{𝐤}{\mathrm{max}}\left\{U\chi ^0(𝐤,\omega =0)\right\},`$ (4)
where $`\chi ^0(𝐪,\omega _l)`$ is given by

$`\chi ^0(𝐪,\omega _l)=T{\displaystyle \underset{𝐤,n}{}}G(𝐪+𝐤,\omega _l+ϵ_n)G(𝐤,ϵ_n)\text{,}`$ (5)
where $`G(𝐪+𝐤,\omega _l+ϵ_n)`$ is the Green function given by the FLEX approximation, and $`\omega _l`$ ($`ϵ_n`$) is the Matsubara frequency for boson (fermion). $`T_\mathrm{N}`$ is determined by the Stoner criterion; $`\alpha _\mathrm{S}=1`$. In the FLEX approximation, however, $`\alpha _\mathrm{S}`$ does not exceed 1 at finite temperatures in two dimensional systems, which is consistent with the Mermin and Wagner theorem. So we determine $`T_\mathrm{N}`$ by the condition $`\alpha _\mathrm{S}=\alpha _\mathrm{N}`$, where we set $`\alpha _\mathrm{N}`$ as $`(1\alpha _\mathrm{N})^1O(100)`$. The AF state will occur through the weak coupling between layers, $`J_{}`$.

Next, We solve the linearized Eliashberg equation to obtain the SC transition temperature $`T_\mathrm{c}`$. For a singlet-pairing case \[$`\varphi (𝐤,ϵ_n)=+\varphi (𝐤,ϵ_n)`$\], it is given by

$`\lambda \varphi (𝐤,ϵ_n)`$ $`=`$ $`T{\displaystyle \underset{𝐪,m}{}}V(𝐤𝐪,ϵ_nϵ_m)`$ (7)
$`\times G(𝐪,ϵ_m)G(𝐪,ϵ_m)\varphi (𝐪,ϵ_m),`$
where $`V(𝐤,\omega _l)=\frac{3}{2}U^2\frac{\chi ^0(𝐤,\omega _l)}{1U\chi ^0(𝐤,\omega _l)}\frac{1}{2}U^2\frac{\chi ^0(𝐤,\omega _l)}{1+U\chi ^0(𝐤,\omega _l)}+U`$. $`T_\mathrm{c}`$ is given by the condition $`\lambda =1`$. In the FLEX approximation, a finite $`T_\mathrm{c}`$ is obtained even in two dimensional systems irrespective of the Hohenberg theorem. However, this approximation gives reasonable $`T_\mathrm{c}`$’s for $`\kappa `$-(ET) organic compounds and high-$`T_\mathrm{c}`$ cuprates . Note that we could not find parameters where the triplet pairing is dominant in the present model.

The obtained phase diagram is given in Fig. 3. 64$`\times `$64 $`k`$-points and 512 Matsubara frequencies are used. The maximum $`T_\mathrm{c}`$ under the condition of $`T_\mathrm{c}>T_\mathrm{N}`$ is about 16K (18K) at $`x=0.8`$ ($`x=0.7`$) for $`U=0.5`$eV ($`U=0.4`$eV). $`T_\mathrm{c}`$ decreases as $`x`$ increases, and $`T_\mathrm{c}6`$K at $`x=1`$ for $`U=0.5`$eV, which is consistent with the observed $`T_\mathrm{c}`$ in several $`\beta `$-(ET)2X compounds. The present study also explains the low $`T_\mathrm{c}`$ in $`\beta `$-(ET)2X at ambient pressure, which corresponds to $`x=1`$ in Fig. 3. On the other hand, we could not find the SC phase in the metallic region for $`x<0`$. We comment that $`T_\mathrm{c}4`$K is obtained at $`x=0`$ for $`U=0.4`$eV as shown in Fig. 3. However, it should be covered with the AF phase because of $`T_\mathrm{N}T_\mathrm{c}`$. Actually, the Stoner factor at 4K is extremely close to 1; $`\alpha _\mathrm{S}0.999`$.

On the other hand, the Néel temperature for $`U=0.5`$eV at $`x=0`$ is 28K (20K) under the condition of $`\alpha _\mathrm{N}=0.99`$ ($`\alpha _\mathrm{N}=0.995`$), which is also consistent with $`T_\mathrm{N}=22`$K in $`\beta ^{}`$-(ET)2ICl2 at ambient pressure. The obtained magnetic order is commensurate \[$`\stackrel{}{Q}=(\pi ,\pi )`$\] for $`x0.7`$ at $`U=0.5`$eV, and it becomes incommensurate for $`x0.8`$ at lower temperatures. The obtained $`T_\mathrm{N}`$ increases as $`x`$ departs from 0, which is interpreted as a natural consequence of the dimensional increase. It is highly demanded to experiment on the pressure dependence of $`T_\mathrm{N}`$. We note that the crossover temperature $`T_{\mathrm{MIT}}`$ determined by the condition $`d\rho /dT=0`$ will be higher than $`T_\mathrm{N}`$ because $`T_{\mathrm{MIT}}`$ reflects the Mott transition, which is beyond the scope of the FLEX approximation.

Figure 4 shows the solution of eq.(7) at $`(k_x,k_y)`$ on the Fermi surface, where $`\theta _𝐤=\mathrm{tan}^1(k_x/k_y)`$. Thus, the obtained SC order parameter is $`d_{x^2\text{-}y^2}`$-wave like, as in the high-$`T_\mathrm{c}`$ cuprates and the $`\kappa `$-(ET) compounds. Note that the present $`x`$-$`y`$ coordinate is different from that of the original crystal; see Fig. 1.

Finally, we study the resistivity $`\rho `$ and the Hall coefficients $`R_\mathrm{H}`$. Considering that the present model lacks the four-fold rotational symmetry, we define $`\rho `$ and $`R_\mathrm{H}`$ as $`\rho =\sqrt{2/(\sigma _{xx}^2+\sigma _{yy}^2)}`$ and $`R_\mathrm{H}=(\sigma _{xy}/B_z)\rho ^2`$ respectively: They are independent of the choice of the $`x`$-$`y`$ coordinate. We calculate them by including the VC’s for the current to maintain the conservation laws . The role of the VC for the current, which is totally dropped in the relaxation time approximation, is very important in strongly correlated systems. For example, the anomalous behaviors of $`R_\mathrm{H}`$ in high-$`T_\mathrm{c}`$ cuprates and in $`\kappa `$-(ET)2X compounds are reproduced only when the VC for the current is taken into account adequately .

Figure 5 shows the obtained results for $`U=0.5`$eV. Experimentally, $`\rho `$ under 8.2GPa continues to increase up to 300K approximately proportional to $`T`$, which is similar to $`\rho `$ in high-$`T_\mathrm{c}`$ cuprates . Such a “bad metal” behavior of $`\rho `$ is well reproduced in the present study. In Fig. 5, $`R_\mathrm{H}`$ for $`x=0.8`$ increases below $`50`$K, which is caused by the VC’s for the current due to the AF fluctuations. However, the obtained enhancement of $`R_\mathrm{H}`$ is rather smaller than that for high-$`T_\mathrm{c}`$ cuprates or $`\beta `$-(ET)2X compounds. As for $`x=0.9`$ and $`1.0`$, the temperature dependence of $`R_\mathrm{H}`$ is much moderate. It is highly demanded to experiment on $`R_\mathrm{H}`$ in $`\beta ^{}`$-(ET)2ICl2. We stress that the observed resistivity under high pressure is considered to be intrinsic free from the lattice contraction effect, which is considerably large in usual organic metals at ambient pressure .

In more detail, however, it might be better to analyze the original tight-binding model in Fig.1, instead of the dimer model, because a transport coefficient is in general sensitive to the shape of the FS.

Here, we discuss the validity of the pressure induced dimensional crossover in $`\beta ^{}`$-(ET)2ICl2. According to Taniguchi, the resistivity along the less-conductive direction (i.e., $`c`$-axis direction) decreases very rapidly as the pressure increases, which suggests the system becomes 2D-like . In fact, the Q1D nature in $`\beta ^{}`$-compound is ascribed to the accidental cancellation of $`t_1`$ in the dimer model although $`|t(q2)|`$ and $`|t(c)|`$ are not so small; see Fig. 1 and Table I. It suggests that the dimensional crossover easily occurs. Moreover, $`\beta ^{}`$-(ET)2AuCl2 remains semiconducting even at 9GPa , which means that the $`\beta ^{}`$-type structure in $`\beta ^{}`$-(ET)2AuCl2 is more robust against the pressure. It is natural because Au-Cl bond length is about 12% smaller than the I-Cl bond length.

As a result, the obtained phase diagram, Fig.3, as well as the concept of the SC driven by the dimensional crossover, will make sense and be reasonable, although the variable $`x`$ in eq.(2) cannot be interpreted as the pressure simply. The large density of states (DOS) at the Fermi level for $`x0.7`$ due to the van Hove singularity around $`(\pm \pi ,0)`$ (see Fig.2) accounts for the high $`T_\mathrm{c}`$ in $`\beta ^{}`$-(ET)2ICl2. Actually, we could not obtain higher $`T_\mathrm{c}`$ ($`\stackrel{>}{}5`$K) under the condition that $`T_\mathrm{c}>T_\mathrm{N}`$ by the slight modification of the original parameters for $`\beta ^{}`$-(ET)2ICl2, as far as the FS is Q1D. We note that $`T_\mathrm{c}`$ in a Q1D system, TMTSF, is very low ($`1`$K), which is recognized by the FLEX approximation .

In the present analysis, we might have underestimated the enhancement of the ratio $`W_{\mathrm{band}}/U`$ ($`W_{\mathrm{band}}`$ being the band width) due to the applied pressure. Actually, the hydrostatic pressure not only distorts the structure ($`\beta ^{}\beta `$), but also shortens the cell volume, which enhances each overlap integral. Then, the enhancement of $`t(p1)`$, which is proportional to $`U_{\mathrm{eff}}`$, is weaker than other hopping parameters .

In summary, we studied the origin of the superconductivity in $`\beta ^{}`$-(ET)2ICl2 based on the FLEX approximation, by assuming that the structure approaches to $`\beta `$-type under the hydrostatic pressure. Our theory predicts that the $`d`$-wave superconductivity occurs in $`\beta ^{}`$-(ET)2ICl2 as a result of the 1D-2D dimensional crossover owing to the crystal structure change under high pressure. Moreover, both $`\rho `$ and $`R_\mathrm{H}`$ were studied in terms of the conserving approximation. Experimental studies on the structure of $`\beta ^{}`$-(ET)2ICl2 as well as the band calculations under the condition of high pressure are highly demanded. In addition, the superconductivity in $`\beta `$-(ET)2X compound is also understood in the same framework.

The author is grateful to H. Taniguchi and M. Miyashita for valuable discussions on experiments. He is also grateful to R. Kato for useful informations on band calculation.