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# Zeros of the Derivatives of Faber Polynomials Associated with a Universal Covering Map

## 1 Introduction

Let $`g`$ be a function which is holomorphic in a neighborhood of infinity such that its Laurent series is of the form

$$g(w)=w+\underset{k=0}{\overset{\mathrm{}}{}}b_kw^k.$$
(1.1)
*Faber polynomials* $`F_k`$ associated with $`g`$ are defined by the generating function

$$\frac{g^{}(w)}{g(w)z}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{F_k(z)}{w^{k+1}},$$
(1.2)
and the *normalized* derivatives of Faber polynomials, $`P_k:=(F_{k+1})^{}/(k+1)`$ for $`k=0,1,2,\mathrm{}`$, satisfy the equation

$$\frac{1}{g(w)z}=\underset{k=0}{\overset{\mathrm{}}{}}\frac{P_k(z)}{w^{k+1}}.$$
(1.3)
For every $`k`$, both $`F_k`$ and $`P_k`$ are monic polynomials of degree $`k`$. To study more about Faber polynomials and their derivatives, see for example ,  and .

Fore each $`k`$, let $`\nu _k`$ be the normalized zero counting measure of $`F_k`$; i.e.,

$$\nu _k=(2\pi k)^1\mathrm{\Delta }(\mathrm{log}|F_k|),$$

where $`\mathrm{\Delta }`$ represents for the generalized Laplacian. Similarly, we denote by $`\mu _k:=(2\pi k)^1\mathrm{\Delta }(\mathrm{log}|P_k|)`$ the normalized zero counting measure of $`P_k`$. Kuijlaars and Saff studied in  the limit behavior of the measures $`\{\nu _k\}`$, and we concern the limit behavior of $`\{\mu _k\}`$. Especially we are interested in the case when $`g`$ is a universal covering map.

Suppose that $`E`$ is a compact set containing more than two points such that $`\mathrm{\Omega }:=\overline{}\backslash E`$ is connected (but not necessarily simply-connected). By the Uniformization theorem (cf. , Chap. 10), there exists a unique number $`r=r(E)>0`$ and a unique normalized universal covering map $`g:\mathrm{\Lambda }(r):=\overline{}\backslash \{|w|r\}\mathrm{\Omega }`$ which has a Laurent expansion of the form (1.1) at infinity. In this case Faber polynomials $`\{F_k\}`$, and their normalized derivatives $`\{P_k\}`$, associated with $`g`$ are also called Faber polynomials, or the normalized derivatives of Faber polynomials, respectively, associated with $`E`$.

Since $`E`$ contains more than two points, the domain $`\mathrm{\Omega }=\overline{}\backslash E`$ carries the unique hyperbolic (Poincaré) metric with constant curvature ($`1`$), and we define $`\stackrel{~}{E}`$ as the union of $`E`$ and the points in $`\mathrm{\Omega }`$ which have more than one shortest curve to $`\mathrm{}`$ with respect to this metric.

###### Example 1.

If $`\overline{}\backslash E`$ is simply-connected, that is, $`E`$ is connected, then $`\stackrel{~}{E}=E`$.

###### Example 2.

Suppose $`E`$ is a compact set consisting of three points. Then $`\stackrel{~}{E}`$ is either a topological tripod or a line segment joining points in $`E`$. See Proposition 18. For example, if $`E=\{1,\eta ,\eta ^2\}`$ where $`\eta =\mathrm{exp}(2\pi i/3)`$ is a third root of unity, then $`\stackrel{~}{E}=\{\eta ^jt:j=0,1,2\text{ and }0t1\}`$ (Figure 1). If $`E=\{1,0,1\}`$, then $`\stackrel{~}{E}=\{z:1z1\}`$ (Figure 2).

There is an alternative way to describe the set $`\stackrel{~}{E}`$. To explain this, let $`\mathrm{\Gamma }`$ be the Fuchsian group such that $`\mathrm{\Lambda }(r)/\mathrm{\Gamma }\mathrm{\Omega }`$ and $`g\tau (z)=g(z)`$ for all $`\tau \mathrm{\Gamma }`$ and $`z\mathrm{\Lambda }(r)`$. We denote by $`\mathrm{\Gamma }_{\mathrm{}}`$ the orbit of $`\mathrm{}`$ under $`\mathrm{\Gamma }`$ and by $`d(,)`$ the hyperbolic distance between two points in $`\mathrm{\Lambda }(r)`$. The set

$$𝒟:=\{w\mathrm{\Lambda }(r):d(\mathrm{},w)<d(\zeta ,w)\text{for all }\zeta \mathrm{\Gamma }_{\mathrm{}},\zeta \mathrm{}\}$$
(1.4)
is called the *Dirichlet domain* associated with $`\mathrm{\Gamma }`$ (or $`g`$) and centered at $`\mathrm{}`$. It is a fundamental region (cf. , Section I-4). Then it will be shown in Corollary 16 that $`\stackrel{~}{E}=g(𝒟)`$, the boundary of the image of $`𝒟`$ under $`g`$.

Let $`𝒰`$ be the collection of compact sets in $``$ with connected complements such that either $`E`$ is not connected, or $`E`$ is connected but $`E`$ contains a singularity other than an outward cusp.

Our main result is:

###### Theorem 3.

1. $`\stackrel{~}{E}`$ is connected.
2. If $`E^{}=\mathrm{}`$, where $`E^{}`$ denotes the interior of $`E`$, then the sequence $`\{\mu _k\}`$ converges in the weak-star topology to the probability measure

$$\mu (z)=\frac{1}{2\pi }\mathrm{\Delta }\mathrm{log}\delta (z)$$
(1.5)
where

$$\delta (z)=\underset{k\mathrm{}}{lim\; sup}|P_k(z)|^{1/k},$$
(1.6)
and the support of $`\mu `$ is $`\stackrel{~}{E}`$. If, in addition, $`\overline{}\backslash E`$ is simply-connected, then $`\mu `$ is the equilibrium distribution of the compact set $`E`$.
3. Suppose $`E^{}`$ is connected and $`E𝒰`$. Then there is a subsequence of $`\{\mu _k\}`$ that converges in the weak-star topology to the measure $`\mu `$ in (1.5), and its support is $`\stackrel{~}{E}`$. If, in addition, $`\overline{}\backslash E`$ is simply-connected, then $`\mu `$ is the equilibrium distribution of the compact set $`E`$.

The polynomials $`\{P_k\}`$ provide some useful tools in other branches of mathematics. For example, A. Atzmon, A. Eremenko and M. Sodin showed in  that if $`a`$ is an element in a complex Banach algebra with unit and $`E`$ is a compact set with connected complement, then the spectrum of $`a`$ is included in $`E`$ if and only if

$$\underset{k\mathrm{}}{lim\; sup}P_k(a)^{1/k}r,$$

where $`\{P_k\}`$ and $`r=r(E)`$ are as before. They also showed in the same paper that for a given analytic germ $`f(z)=_{k=0}^{\mathrm{}}f_kz^{k1}`$ at infinity with values in a Banach space, the polynomials $`\{P_k\}`$ can be used to determine whether $`f`$ has analytic continuation to $`\mathrm{\Omega }=\overline{}\backslash E`$ or not.

For a formal Laurent series of the form (1.1), it is known (cf. ) that the zeros of $`P_k(z)`$ are exactly the eigenvalues of the leading $`k\times k`$ principal submatrix of the infinite Toeplitz matrix

$$\left[\begin{array}{ccccc}b_0& b_1& b_2& b_3& \mathrm{}\\ 1& b_0& b_1& b_2& \mathrm{}\\ 0& 1& b_0& b_1& \mathrm{}\\ 0& 0& 1& b_0& \mathrm{}\\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}\end{array}\right].$$

## 2 Known facts

Suppose the Laurent expansion of $`g`$ is given by (1.1), and we assume that $`\rho _0:=lim\; sup_k\mathrm{}|b_k|^{1/k}<\mathrm{}`$ so that the series is convergent in $`|w|>\rho _0`$. We also define $`\delta _0`$ as the smallest nonnegative number such that $`g`$ has a *meromorphic* extension to $`\mathrm{\Lambda }(\delta _0)=\{|w|>\delta _0\}\{\mathrm{}\}`$. Now for every $`z`$, we write $`g^1(z):=\{w\mathrm{\Lambda }(\rho _0):g(w)=z\}`$ and $`\stackrel{~}{g}^1(z):=\{w\mathrm{\Lambda }(\delta _0):g(w)=z\}`$.

Note that the multiplicity of a point $`w\mathrm{\Lambda }(\delta _0)`$ is $`m`$ if $`g^{}(w)=\mathrm{}=g^{(m1)}(w)=0`$ and $`g^{(m)}(w)0`$. The following definition is due to J. L. Ullman (, ).

###### Definition 4.

For every nonnegative integer $`p`$, $`\stackrel{~}{C}_p`$ (or $`C_p`$) is the set of all points $`z`$ such that the points of largest absolute value in $`\stackrel{~}{g}^1(z)`$ (or $`g^1(z)`$, respectively) have total multiplicity $`p`$.

From the definition, one can easily see that $`\stackrel{~}{C}_0=\backslash g(\mathrm{\Lambda }(\delta _0))`$ is compact and $`\stackrel{~}{C}_1`$ is an open set containing a neighborhood of infinity.

The next two statements are analogous to Theorem 1.3 and Theorem 1.4 of , where the theorems below are proved with $`C_0,C_1,F_k`$ and $`\nu _k`$ in place of $`\stackrel{~}{C}_0,\stackrel{~}{C}_1,P_k`$ and $`\mu _k`$, respectively.

###### Theorem 5 (cf. , Theorem 1.3).

If the interior of $`\stackrel{~}{C}_0`$ is empty, then the sequence $`\{\mu _k\}`$ converges in the weak-star topology to the measure $`\mu `$ in (1.5) and the support of $`\mu `$ is equal to $`\stackrel{~}{C}_1`$. If, in addition, $`\overline{}=\stackrel{~}{C}_0\stackrel{~}{C}_1`$, then $`\mu `$ is the equilibrium distribution of the compact set $`\stackrel{~}{C}_0`$.

###### Theorem 6 (cf. , Theorem 1.4).

If the interior of $`\stackrel{~}{C}_0`$ is connected, then there is a subsequence of $`\{\mu _k\}`$ that converges in the weak-star topology to the measure $`\mu `$ in (1.5), and the support of $`\mu `$ is $`\stackrel{~}{C}_1`$. If, in addition, $`\overline{}=\stackrel{~}{C}_0\stackrel{~}{C}_1`$, then $`\mu `$ is the equilibrium distribution of the compact set $`\stackrel{~}{C}_0`$.

The proofs of Theorems 5 and 6 are given in Appendix A. In fact, one may check that these theorems can be shown by the same arguments for Theorems 1.3 and 1.4 in , provided that Ullman’s results used in  are replaced by those listed below.

###### Lemma 7 (, Lemma 7; cf. , Lemma 2.2).

Every $`z_0\stackrel{~}{C}_p`$, $`p2`$, has a neighborhood $`B(z_0,ϵ):=\{z:|zz_0|<ϵ\}`$ such that $`\stackrel{~}{C}_1\overline{B(z_0,ϵ)}`$ consists of a finite number of analytic Jordan arcs each joining $`z_0`$ to a point on the circle $`B(z_0,ϵ)`$. Any two arcs intersect only at $`z_0`$. The remaining points of $`B(z_0,ϵ)`$ are in $`\stackrel{~}{C}_1`$.

From Lemma 7, we have an important corollary.

###### Corollary 8 (cf. , Corollary 2.3).

$$\stackrel{~}{C}_1=\stackrel{~}{C}_0\underset{p2}{}\stackrel{~}{C}_p.$$

A point $`\xi `$ is called a limit point of the zeros of $`\{P_k\}`$ if there exist an increasing sequence $`\{k_j\}`$ and a zero $`\xi _j`$ of $`P_{k_j}`$ for each $`j`$ such that $`\xi =lim_j\mathrm{}\xi _j`$.

###### Theorem 9 (; cf. , Theorem 1.2).

All limit points of the zeros of $`\{P_k\}`$ are in $`\backslash \stackrel{~}{C}_1`$. Every boundary point of $`\stackrel{~}{C}_1`$ is a limit point of the zeros of $`\{P_k\}`$.

###### Proof.

The first statement is Lemma 3 of , and the second statement is given in the proof of Lemma 11 in the same paper. ∎

###### Lemma 10 (, Lemma 12).

The boundary of the unbounded component of $`\stackrel{~}{C}_1`$ is a connected set.

###### Lemma 11 (, Lemma 8; cf. , Lemma 2.4).

For every $`ϵ_0>0`$ and $`z_0\stackrel{~}{C}_p`$, $`p2`$, there are $`ϵ_1>0`$ and $`z_1\stackrel{~}{C}_q`$, $`q2`$, such that

$$B(z_1,ϵ_1)B(z_0,ϵ_0),$$

$$B(z_1,ϵ_1)\stackrel{~}{C}_0=\mathrm{},$$

$$B(z_1,ϵ_1)\stackrel{~}{C}_1=D_1D_2,$$

where $`D_1`$ and $`D_2`$ are disjoint non-empty domains. Moreover, there exist analytic functions $`f_1`$ and $`f_2`$ on $`B(z_1,ϵ_1)`$ such that

$`|f_1(z)|>|f_2(z)|,`$ $`zD_1,`$ (2.1)
$`|f_i(z)|=\delta (z),`$ $`zD_i,i=1,2,`$ (2.2)
where $`\delta (z)`$ is defined in (1.6).

###### Lemma 12 (; cf. , Lemma 3.1).

1. For every $`z\stackrel{~}{C}_0`$,

$$\delta (z)=\delta _0,$$
(2.3)
where $`\delta _0`$ is defined in the first paragraph of this section.
2. For every $`z\stackrel{~}{C}_0`$,

$$\delta (z)=\mathrm{max}\{|w|:w\stackrel{~}{g}^1(z)\}.$$
(2.4)
3. For every $`z\stackrel{~}{C}_1`$,

$$\delta (z)=\underset{k\mathrm{}}{lim}|P_k(z)|^{1/k};$$
(2.5)
i.e., the $`lim\; sup`$ in (1.6) can be replaced by $`lim`$.

An immediate corollary of (2.3) and (2.4) is:

###### Corollary 13 (cf. , Lemma 3.1).

$`\delta (z)`$ is a continuous function on $``$.

## 3 Proof of Theorem 3

Suppose $`E`$ is a compact set containing more than two points such that its complement $`\mathrm{\Omega }=\overline{}\backslash E`$ is connected, and let $`g:\mathrm{\Lambda }(r)\mathrm{\Omega }`$ be the uniformizing map with Laurent expansion (1.1) at infinity.

###### Lemma 14.

For given $`z\mathrm{\Omega }`$, there is a one-to-one correspondence between shortest curves (with respect to the hyperbolic metric in $`\mathrm{\Omega }`$) from $`z`$ to $`\mathrm{}`$ and points in $`\stackrel{~}{g}^1(z)`$ with largest absolute value.

###### Proof.

Suppose $`\gamma :[a,b]\mathrm{\Omega }`$ is a shortest curve from $`z`$ to infinity such that $`\gamma (a)=z`$ and $`\gamma (b)=\mathrm{}`$, and let $`\stackrel{~}{\gamma }:[a,b]\mathrm{\Lambda }(r)`$ be the lifting curve of $`\gamma `$ such that $`\stackrel{~}{\gamma }(b)=\mathrm{}`$. Because $`g`$ is a local isometry between $`\mathrm{\Lambda }(r)`$ and $`\mathrm{\Omega }`$ with hyperbolic metrics on them, we know that $`\stackrel{~}{\gamma }`$ is a shortest curve in $`\mathrm{\Lambda }(r)`$; i.e., the trace of $`\stackrel{~}{\gamma }`$ is the ray

$$\{tw:1t\mathrm{}\}$$
(3.1)
for $`w=\stackrel{~}{\gamma }(a)`$. Now if $`|w_0|>|\stackrel{~}{\gamma }(a)|`$ for some $`w_0\stackrel{~}{g}^1(z)`$, let $`\alpha `$ be the ray (3.1) with $`w=w_0`$. Then since $`|w_0|>|\stackrel{~}{\gamma }(a)|`$, $`\alpha `$ is shorter than $`\stackrel{~}{\gamma }`$, hence $`g(\alpha )`$ is shorter than $`g(\stackrel{~}{\gamma })=\gamma `$. Since $`g(\alpha )`$ connects $`z`$ to $`\mathrm{}`$, this contradicts to the choice of $`\gamma `$, proving $`|\stackrel{~}{\gamma }(a)|=\mathrm{max}\{|w|:w\stackrel{~}{g}^1(z)\}`$. Conversely, if we have a point $`w_0\stackrel{~}{g}^1(z)`$ such that $`|w_0|=\mathrm{max}\{|w|:w\stackrel{~}{g}^1(z)\}`$, then by the same argument above the image of the ray (3.1) with $`w=w_0`$ is a shortest curve in $`\mathrm{\Omega }`$ connecting $`z=g(w_0)`$ and $`\mathrm{}`$. Therefore the map $`\gamma \stackrel{~}{\gamma }(a)`$ is bijective between the set of shortest curves from $`z`$ to $`\mathrm{}`$ and the set of points in $`\stackrel{~}{g}^1(z)`$ with largest absolute value. The lemma follows. ∎

Note that $`g^{}(w)0`$ for any $`w\mathrm{\Lambda }(r)`$ since $`g`$ is a covering map. Therefore the following corollary is an immediate consequence of Lemma 14.

###### Corollary 15.

For each $`z\mathrm{\Omega }`$, $`z\stackrel{~}{E}`$ if and only if $`z\stackrel{~}{C}_p`$ for some $`p2`$.

Let $`𝒟`$ be the Dirichlet domain associated with $`g`$ and centered at $`\mathrm{}`$, as we introduced in Section 1.

###### Corollary 16.

$`\stackrel{~}{E}=g(𝒟)`$.

###### Proof.

Because $`𝒟`$ is a fundamental region, $`g(\overline{𝒟})=\mathrm{\Omega }`$ and $`g|_𝒟`$ is injective. Therefore $`E=\mathrm{\Omega }g(𝒟)`$, $`E^{}g(𝒟)=\mathrm{}`$, and $`g(𝒟)\mathrm{\Omega }=g(𝒟)`$. Hence to prove the corollary, it suffices to show that $`\stackrel{~}{E}\mathrm{\Omega }=g(𝒟)\mathrm{\Omega }=g(𝒟)`$.

Let $`d(,)`$ denote the hyperbolic distance between two points in $`\mathrm{\Lambda }(r)`$, and let $`\mathrm{\Gamma }`$ be the corresponding Fuchsian group such that $`\mathrm{\Lambda }(r)/\mathrm{\Gamma }\mathrm{\Omega }`$. Then for each $`w_1,w_2\mathrm{\Lambda }(r)`$ and $`\tau \mathrm{\Gamma }`$, $`d(w_1,w_2)=d(\tau (w_1),\tau (w_2))`$. Therefore, $`d(w,\mathrm{})d(w,\tau (\mathrm{}))`$ for all $`\tau \mathrm{\Gamma }`$ if any only if $`d(w,\mathrm{})d(\tau (w),\mathrm{})`$ for all $`\tau \mathrm{\Gamma }`$, or $`|w||\tau (w)|`$ for all $`\tau \mathrm{\Gamma }`$. This means that $`w\overline{𝒟}`$ if and only if $`|w|=\mathrm{max}\{|\xi |:\xi \stackrel{~}{g}^1(g(w))\}`$. Now the corollary follows from Lemma 14 because we have $`w𝒟`$ if and only if there exist $`w^{}𝒟\backslash \{w\}`$ and $`\tau \mathrm{\Gamma }`$ such that $`w=\tau (w^{})`$, that is, $`g(w)=g(w^{})`$ (, p. 37). ∎

Note that the proof of this corollary also shows that $`\mathrm{\Omega }\backslash \stackrel{~}{E}=g(𝒟)`$.

###### Proof of Theorem 3.

Suppose $`E𝒰`$. Then $`\mathrm{\Omega }`$ is simply-connected, hence $`\stackrel{~}{E}=E=\mathrm{\Omega }`$ is connected. If $`E𝒰`$, then as shown in  (p. 444), $`g`$ cannot be extended to $`\mathrm{\Lambda }(r_0)`$ for any $`r_0<r`$; i.e., $`r=\delta _0`$, where $`\delta _0`$ is defined in the first paragraph of Section 2. Therefore we have $`E=\stackrel{~}{C}_0`$, hence $`\mathrm{\Omega }\backslash \stackrel{~}{E}=\stackrel{~}{C}_1`$ by Corollary 15; i.e., $`\stackrel{~}{C}_1=g(𝒟)`$. Now Lemma 10 implies the statement (1) of Theorem 3, since $`\stackrel{~}{C}_1=g(𝒟)`$ is connected and $`\stackrel{~}{E}=g(𝒟)`$.

To prove (2) of Theorem 3, assume that $`E^{}=\mathrm{}`$. If $`E`$ is not connected, then we have $`\delta _0=r`$, or $`\stackrel{~}{C}_0=E`$, hence $`\stackrel{~}{C}_1=\stackrel{~}{E}`$ by Corollaries 8 and 15. Since $`\stackrel{~}{C}_0=E`$ has empty interior, Theorem 5 implies that $`\{\mu _k\}`$ converges to $`\mu `$ and its support is $`\stackrel{~}{C}_1=\stackrel{~}{E}`$. If $`E`$ is connected, we consider two cases: $`\delta _0<r`$ and $`\delta =r`$. If $`\delta _0<r`$, then every point $`zE`$ corresponds to at least 2 points on $`|w|=r`$ counting muliplicity, hence $`\stackrel{~}{C}_0=\mathrm{}`$, $`\delta (z)`$ is constant ($`=r`$) on $`E`$, and $`E=_{p2}\stackrel{~}{C}_p=\stackrel{~}{C}_1`$. If $`\delta _0=r`$, then $`\delta (z)`$ is constant ($`=r=\delta _0`$) on $`E=\stackrel{~}{C}_0`$ by Lemma 12(a) and $`E=\stackrel{~}{E}=\stackrel{~}{C}_1`$ as before. In any case $`E=\stackrel{~}{E}=\stackrel{~}{C}_1`$ and $`\delta (z)`$ is constant on $`E`$. Therefore by Theorem 5, the sequence $`\{\mu _k\}`$ converges to $`\mu =(2\pi )^1\mathrm{\Delta }\mathrm{log}\delta `$ and its support is $`E=\stackrel{~}{E}`$. Furthermore, it is the equilibrium distribution of $`E`$ since $`\delta `$ is constant on $`E`$. This proves (2) of Theorem 3.

The statement (3) of Theorem 3 follows from Theorem 6 by the same argument as above. ∎

One cannot replace $`\mu _k`$ in Theorem 3 by $`\nu _k=(2\pi k)^1\mathrm{\Delta }(\mathrm{log}|F_k|)`$. To see this, let $`\eta =\mathrm{exp}(2\pi i/3)`$ , $`E=\{1,\eta ,\eta ^2\}`$, and $`g:\mathrm{\Lambda }(r)\mathrm{\Omega }`$ the universal covering map with Laurent expansion (1.1). As in Example 2,

$$\stackrel{~}{E}=\{\eta ^jt:j=0,1,2\text{ and }0t1\}.$$

Now we claim that there is a subsequence of $`\{\nu _k\}`$ such that the support of its weak-star limit $`\nu `$ is different from $`\stackrel{~}{E}`$ (Figure 4).

Let $`\rho _0`$ be the maximum of absolute values of the finite poles of $`g`$, and let $`\mathrm{\Omega }_0:=g(\mathrm{\Lambda }(\rho _0))`$. Then since $`\rho _0>r=r(E)`$ and

$$\eta ^1g(\eta w)=g(w)\text{for all }w\mathrm{\Lambda }(r),$$
(3.2)
one can easily see that $`\backslash \mathrm{\Omega }_0=𝒩_0𝒩_1𝒩_2`$, where $`𝒩_0`$ is a closed neighborhood of 1 and $`𝒩_j=\eta ^j𝒩_0`$, $`j=1,2`$. Therefore from Definition 4 and Corollary 2.3 of  (cf. Corollary 8),

$$C_1=C_0\underset{p2}{}C_p=\underset{j=1}{\overset{3}{}}𝒩_j\left(\stackrel{~}{E}\mathrm{\Omega }_0\right),$$

which is the set shown in Figure 4.

By Lemma 3.1 of  (cf. Lemma 12), we have $`lim\; sup_k\mathrm{}|F_k(1)|=\rho _0`$. Thus there exists a subsequence $`\{F_{k_j}\}`$ of $`\{F_k\}`$ such that $`lim_j\mathrm{}|F_{k_j}(1)|=\rho _0`$. Considering (3.2), we also have $`lim_j\mathrm{}|F_{k_j}(\eta )|=lim_j\mathrm{}|F_{k_j}(\eta ^2)|=\rho _0`$. Now by Lemma 5.1 and the proof of Lemma 5.2 of  (cf. Lemmas 20 and 21), the sequence $`\{\nu _{k_j}\}`$ converges in the weak-star topology to the measure

$$\nu (z)=\frac{1}{2\pi }\mathrm{\Delta }\mathrm{log}\left(\underset{k\mathrm{}}{lim\; sup}|F_k(z)|^{1/k}\right),$$

and by Lemma 4.1 of  (cf. Lemma 19) the support of $`\nu `$ is $`C_1`$.

## 4 Compact sets consisting of three points

Let $`\mathrm{\Omega }_{0,1}:=\backslash \{0,1\}`$ be a metric space equipped with the hyperbolic metric on it, and we use the notations $`H^+:=\{Im(z)>0\}`$, $`H^{}:=\{Im(z)<0\}`$, $`I_1:=\{z:z<0\},`$ $`I_2:=\{z:0<z<1\},`$ $`I_3:=\{z:z>1\}`$, and $`I:=I_1I_2I_3`$. The following lemma is very trivial but plays a crucial role in our arguments later.

###### Lemma 17.

Suppose $`z_1H^+`$ and $`z_2H^+I`$. Then there is a unique shortest curve $`\gamma `$ connecting $`z_1`$ and $`z_2`$. Moreover, the interior arc of $`\gamma `$ does not intersect $`I`$.

###### Proof.

The proof is omitted here and left to the reader. In fact, one may prove it by considering the symmetric property of $`\mathrm{\Omega }_{0,1}`$. ∎

Let $`𝔻:=\{z:|z|<1\}`$ with the hyperbolic metric on it, and we assume that $`G:𝔻\mathrm{\Omega }_{0,1}`$ is a holomorphic universal covering map such that for a given point $`aI_1H^+`$, $`G(0)=a`$. Let $`\mathrm{\Gamma }`$ be the modular group on $`𝔻`$ such that $`𝔻/\mathrm{\Gamma }\mathrm{\Omega }_{0,1}`$ and $`G\tau (z)=G(z)`$ for all $`\tau \mathrm{\Gamma }`$, and we denote by $`\mathrm{\Gamma }_0`$ the orbit of the origin under $`\mathrm{\Gamma }`$ and by $`𝒟_0`$ the Dirichlet domain with centered at the origin. Note that each component of $`G^1(H^+)`$ or $`G^1(H^{})`$ is a hyperbolic open triangle, and each side of such a triangle is a geodesic curve and a component of $`G^1(I_j)`$ for some $`j=1,2,3`$. Now we consider the cases $`aI_1`$ and $`aH^+`$ separately.

If $`aI_1`$, there are two hyperbolic triangles $`\mathrm{}^+`$ and $`\mathrm{}^{}`$ in $`𝔻`$ such that $`G(\mathrm{}^+)=H^+`$, $`G(\mathrm{}^{})=H^{}`$ and $`0\overline{\mathrm{}^+}\overline{\mathrm{}^{}}`$. Suppose $`w\mathrm{}^+`$. Then by Lemma 17, there exists a unique shortest curve $`\gamma :[t_0,t_1]H^+\{a\}`$ such that $`\gamma (t_0)=a`$ and $`\gamma (t_1)=G(w)`$. Let $`\stackrel{~}{\gamma }`$ be the lifting curve such that $`\stackrel{~}{\gamma }(t_0)=0`$. Because $`\gamma `$ does not intersect $`I`$, $`\stackrel{~}{\gamma }`$ cannot intersect $`G^1(I)`$, hence $`\stackrel{~}{\gamma }(t_1)=w`$. Now if $`G(w^{})=G(w)`$ for some $`w^{}w`$, the shortest curve $`\alpha `$ from $`w^{}`$ to $`0`$ intersects $`G^1(I)`$, thus $`G(\alpha )`$ intersects $`I`$. Therefore the length of $`G(\alpha )`$ is strictly greater than the length of $`\gamma `$, or the the length of $`\alpha `$ is strictly greater than the length of $`\stackrel{~}{\gamma }`$. This shows that $`w𝒟_0`$. Because a similar argument holds for any $`w\mathrm{}^{}`$, we have

$$\mathrm{}^+\mathrm{}^{}𝒟_0.$$
(4.1)
Since $`𝒟_0`$ is a fundamental region, $`G`$ is univalent in $`𝒟_0`$ and $`G(\overline{𝒟}_0)=\mathrm{\Omega }_{0,1}`$. Thus (4.1) in fact shows that

$$𝒟_0=\left(\overline{\mathrm{}^+}\overline{\mathrm{}^{}}\right)^{}.$$

Note that in this case $`G(𝒟_0)=I_2I_3`$.

We next consider the case $`aH^+`$. Let $`\mathrm{}_0`$ be the triangle such that $`0\mathrm{}_0`$ and $`G(\mathrm{}_0)=H^+`$. For each $`j\{1,2,3\}`$, we denote by $`L_j`$ the side of $`\mathrm{}_0`$ such that $`G(L_j)=I_j`$, and let $`\mathrm{}_j`$ be the hyperbolic triangle which is obtained by reflecting $`\mathrm{}_0`$ with respect to $`L_j`$. Similarly, we denote by $`\mathrm{}_{j,k}`$ the triangle obtained by reflecting $`\mathrm{}_j`$ through the side over $`I_k`$, $`kj`$. Finally, let $`\zeta _{j,k}`$ be the point in $`\mathrm{}_{j,k}`$ such that $`G(\zeta _{j,k})=a`$. See Figure 5.

By Lemma 17 and the same argument as above, $`\overline{\mathrm{}}_0𝒟_0`$. Similarly, the closed triangle $`\overline{\mathrm{}}_{j,k}`$ is contained in the Dirichlet domain with center at $`\zeta _{j,k}`$. Since $`𝒟_0`$ is connected, we have

$$\overline{\mathrm{}}_0𝒟_0\underset{j=0}{\overset{4}{}}\mathrm{}_j\underset{j=1}{\overset{3}{}}L_j.$$
(4.2)
Let $`A:=\{\zeta _{j,k}:1j,k3,jk\}`$. Then for all $`w𝒟_0`$ and $`\zeta \mathrm{\Gamma }_0\backslash (\{0\}A)`$, the shortest curve connecting these two points must pass through a point $`w^{}\mathrm{}_{j,k}`$ for some $`j`$ and $`k`$, $`jk`$. Therefore, denoting by $`d_𝔻`$ the hyperbolic distance in $`𝔻`$, we have

$$d_𝔻(w,\zeta )=d_𝔻(w,w^{})+d_𝔻(w^{},\zeta )>d_𝔻(w,w^{})+d_𝔻(w^{},\zeta _{j,k})d_𝔻(w,\zeta _{j,k}),$$

because $`w^{}\mathrm{}_{j,k}`$ is contained in the Dirichlet domain with center at $`\zeta _{j,k}`$. This implies that

$$𝒟_0=\underset{\zeta A}{}\{w:d_𝔻(w,0)<d_𝔻(w,\zeta )\}.$$
(4.3)
We obtained $`\mathrm{}_{j,k}`$ by reflecting $`\mathrm{}_0`$ through a side over $`I_j`$ and then a side over $`I_k`$. The triangle $`\mathrm{}_0`$ can be obtained from $`\mathrm{}_{k,j}`$ by the same way. Therefore, there is $`\tau \mathrm{\Gamma }`$ such that $`\tau (0)=\zeta _{j,k}`$ and $`\tau (\zeta _{k,j})=0`$. Then with the notation

$$\mathrm{}_{j,k}:=\{w:d_𝔻(w,0)=d_𝔻(w,\zeta _{j,k})\},$$

we have $`\tau (\mathrm{}_{k,j})=\mathrm{}_{j,k}`$; i.e., $`G(\mathrm{}_{j,k})=G(\mathrm{}_{k,j})`$. Also note that $`\mathrm{}_{j,k}`$ is a geodesic curve which separates the two sides of $`\mathrm{}_j`$ lying over $`I_j`$ and $`I_k`$, because $`\overline{\mathrm{}}_0`$ and $`\overline{\mathrm{}}_{j,k}`$ are contained in the Dirichlet domains with centers at $`0`$ and $`\zeta _{j,k}`$, respectively. Since each sides of $`\mathrm{}_j`$ are also geodesic, we conclude that one end of $`\mathrm{}_{j,k}`$ approaches to the common vertex (at infinity) of $`\mathrm{}_0`$ and $`\mathrm{}_{j,k}`$, and $`\mathrm{}_{j,k}`$ intersects the side of $`\mathrm{}_j`$ which is over $`I_l`$, $`lj,k`$. In particular, this implies that $`\mathrm{}_{j,k}\mathrm{}_{j,k^{}}\mathrm{}`$ for $`kk^{}`$.

Now one can easily see that $`𝒟_0`$ is a hexagon with exactly 3 vertices at infinity, which are in fact the vertices of $`\mathrm{}_0`$. Moreover along $`𝒟_0`$, finite and infinite vertices are placed alternatively and the $`G`$-images of the two sides sharing a common infinite vertex are same. Therefore, the three finite vertices are mapped to the same point, say $`b`$, and $`G(𝒟_0)`$ is a tripod with center at $`b`$ such that each leg of it is a hyperbolic geodesic curve connecting $`b`$ to one of the points $`0,1,\mathrm{}`$. Now we are ready to prove the following proposition.

###### Proposition 18.

Suppose $`E`$ consists of three points. If the points in $`E`$ are on a straight line, $`\stackrel{~}{E}`$ is the line segment connecting points in $`E`$; i.e., $`\stackrel{~}{E}`$ is the convex hull of $`E`$. Otherwise $`\stackrel{~}{E}`$ is a tripod.

###### Proof.

Let $`g:\mathrm{\Lambda }(r)\overline{}\backslash E`$ be a holomorphic universal map. We choose a linear transformation $`T`$ such that $`T(E)=\{0,1,\mathrm{}\}`$ and $`a:=T(\mathrm{})H^+I_1`$. Since every linear transformation maps circles onto circles, we see that $`aI_1`$ if and only if $`\mathrm{}`$ is on the circle passing through the points in $`E`$; i.e., if and only if the points in $`E`$ are on a straight line. Now let $`G(w):=Tg(r/w)`$. Then $`G`$ is a holomorphic universal covering map from $`𝔻`$ to $`\mathrm{\Omega }_{0,1}`$ such that $`G(0)=a`$. Now the proposition follows from Corollary 16 and the arguments preceding the proposition, since $`T`$ is a conformal map and the map $`wr/w`$ sends the Dirichlet domain in $`𝔻`$ with center at the origin onto the Dirichlet domain in $`\mathrm{\Lambda }(r)`$ with center at infinity. ∎

## Appendix A Appendix: Proofs of Theorems 5 and 6

For a measure $`\omega `$ on $``$, we denote its logarithmic potential by

$$𝒫_\omega (z):=_{}\mathrm{log}|z\xi |d\omega (\xi )$$

and let $`𝒫(z):=\mathrm{log}\delta (z)=lim\; sup_k\mathrm{}k^1\mathrm{log}|P_k(z)|`$. Note that $`\mu (z)=(2\pi )^1\mathrm{\Delta }𝒫(z)`$ by the definition of $`𝒫`$ and (1.5).

###### Lemma 19.

1. $`𝒫`$ is subharmonic on $``$.
2. $`𝒫`$ is harmonic on $`\stackrel{~}{C}_1(\stackrel{~}{C}_0)^{}`$, but not at points of $`\stackrel{~}{C}_1`$.
3. $`\mu `$ is a probability measure with support $`\stackrel{~}{C}_1`$.
4. $`𝒫_\mu (z)=𝒫(z)`$ for all $`z`$.

###### Proof.

Note that the limit superior of a sequence of subharmonic functions is subharmonic if it is upper semi-continuous. Since $`𝒫`$ is upper semi-continuous by Corollary 13 and $`k^1\mathrm{log}|P_k|`$ is subharmonic for all $`k`$, (a) follows.

By (2.3), $`𝒫`$ is constant on $`\stackrel{~}{C}_0`$ hence harmonic on $`(\stackrel{~}{C}_0)^{}`$. If $`z\stackrel{~}{C}_1`$, there exist a neighborhood $`N`$ of $`z`$ and an inverse branch of $`g`$, say $`f`$, defined on $`N`$ such that

$$|f(\xi )|=\mathrm{max}\{|\zeta |:\zeta \stackrel{~}{g}^1(\xi )\},\text{for }\xi N.$$

Therefore by (2.4), $`𝒫(\xi )=\mathrm{log}|f(\xi )|`$ in $`N`$ hence harmonic since $`|f(\xi )|>\delta _00`$.

The function $`𝒫`$ is not harmonic on $`\stackrel{~}{C}_0`$ since Lemma 12 implies that $`𝒫(z)=\delta _0`$ for all $`z\stackrel{~}{C}_0`$ but $`𝒫(z)>\delta _0`$ if $`z\stackrel{~}{C}_0`$. We next show that $`𝒫`$ is not harmonic at a point $`z\stackrel{~}{C}_p`$, $`p2`$. If it is not the case, there exists a neighborhood $`N`$ of $`z`$ such that $`𝒫`$ is harmonic on $`N`$. Then by Lemma 11, there exist a subdomain $`N^{}N`$, two disjoint domains $`D_1,D_2`$ such that $`D_1D_2=N^{}\stackrel{~}{C}_1`$, and analytic functions $`f_1`$ and $`f_2`$ satisfying (2.1) and (2.2). This implies $`𝒫(\xi )\mathrm{log}|f_2(\xi )|`$ is a harmonic function which is positive on $`D_1`$ and zero on $`D_2`$, which is impossible. Therefore $`𝒫`$ is not harmonic on $`_{p2}\stackrel{~}{C}_p`$. Now (b) follows from Corollary 8.

Since $`𝒫`$ is subharmonic by (a), $`\mu =(2\pi )^1\mathrm{\Delta }𝒫`$ is a measure. Moreover by (b), the support of $`\mu `$ is $`\stackrel{~}{C}_1`$. Therefore to show (c), it suffices to show that $`\mu ()=1`$. Note that since $`\mu `$ has compact support and $`𝒫`$ is harmonic off the set $`\stackrel{~}{C}_1`$, the Riesz Decomposition Theorem (, Theorem II.21) implies that

$$u(z):=𝒫_\mu (z)+𝒫(z)$$
(A.1)
is harmonic on any bounded domain $`D`$ containing $`\stackrel{~}{C}_1`$. Because this is true for any arbitrary large domain $`D`$, $`u`$ is in fact harmonic on $``$, hence constant.

Let $`f`$ be the inverse of $`g`$ defined on a neighborhood of $`\mathrm{}`$ such that $`f(\mathrm{})=\mathrm{}`$. Then it is easy to see from (1.1) that $`f(z)=z+O(1)`$ as $`z\mathrm{}`$. Moreover by (2.4), $`\delta (z)=|f(z)|`$ for sufficiently large $`z`$. Thus $`𝒫(z)=\mathrm{log}|f(z)|=\mathrm{log}|z|+o(1)`$ as $`z\mathrm{}`$. But since $`𝒫_\omega (z)=\omega ()\mathrm{log}|z|+o(1)`$ for any finite measure $`\omega `$ with compact support, we conclude from (A.1) that $`\mu ()=1`$ and $`u0`$, which shows (c) and (d) simultaneously. This completes the proof. ∎

Recall that $`\mu _k(z)=(2\pi k)^1\mathrm{\Delta }(\mathrm{log}|P_k(z)|)`$. Since $`P_k`$ is a monic polynomial of degree $`k`$, it can be shown that $`𝒫_{\mu _k}=k^1\mathrm{log}|P_k|`$ by the same argument for Lemma 19(d).

###### Lemma 20.

Let $`\omega `$ be any weak-star limit of $`\{\mu _k\}`$. Then

$`𝒫_\omega (z)=𝒫_\mu (z),`$ $`z\stackrel{~}{C}_1\stackrel{~}{C}_1,`$ (A.2)
$`𝒫_\omega (z)𝒫_\mu (z),`$ $`z.`$ (A.3)
###### Proof.

Suppose $`\{\mu _{k_j}\}`$ converges to $`\omega `$ in the weak-star topology. By (2.5) and Lemma 19(d),

$$\underset{j\mathrm{}}{lim}𝒫_{\mu _{k_j}}(z)=\underset{j\mathrm{}}{lim}k_j^1\mathrm{log}|P_{k_j}(z)|=\mathrm{log}\delta (z)=𝒫(z)=𝒫_\mu (z)$$

for all $`z\stackrel{~}{C}_1`$. Therefore the Lower Envelope Theorem (, Theorem 3.8) implies that $`𝒫_\mu (z)=𝒫_\omega (z)`$ for all $`z\stackrel{~}{C}_1`$ except on a set of logarithmic capacity zero. On the other hand, Theorem 9 implies that the support of $`𝒫_\omega `$ is contained in $`\backslash \stackrel{~}{C}_1`$, hence it is harmonic in $`\stackrel{~}{C}_1`$. Because $`𝒫_\mu =𝒫(z)`$ is also harmonic in $`\stackrel{~}{C}_1`$ (Lemma 19(b)), we conclude that $`𝒫_\mu (z)=𝒫_\omega (z)`$ for all $`z\stackrel{~}{C}_1`$.

Corollary 13 implies that $`𝒫_\mu (z)=\mathrm{log}\delta (z)`$ is continuous if $`\delta (z)>0`$. If $`\delta (z)=0`$ (this happens only when $`z\stackrel{~}{C}_0`$ and $`\delta _0=0`$), $`𝒫_\mu (z)=\mathrm{}`$. Therefore we have

$$𝒫_\omega (z)𝒫_\mu (z),\text{for all }z\stackrel{~}{C}_1,$$
(A.4)
because $`𝒫_\omega `$ is lower semi-continuous and $`𝒫_\omega =𝒫_\mu `$ on $`\stackrel{~}{C}_1`$.

If $`z\stackrel{~}{C}_1`$ and $`z`$ approaches to $`\stackrel{~}{C}_0`$, we have $`𝒫_\omega (z)=𝒫_\mu (z)\mathrm{log}\delta _0`$. Therefore, the minimum principle implies that $`𝒫_\omega (z)\mathrm{log}\delta _0=𝒫_\mu (z)`$ for all $`z\stackrel{~}{C}_0`$. Note that combining this result with (A.4), we also have $`𝒫_\omega (z)=𝒫_\mu (z)`$ for all $`z\stackrel{~}{C}_0`$.

Now it remains to show that $`𝒫_\omega (z)𝒫_\mu (z)`$ for all $`z\stackrel{~}{C}_p`$, $`p2`$. But for sufficiently small $`ϵ`$, Lemma 7 implies that the circle $`\{\xi :|z\xi |=ϵ\}`$ is contained in $`\stackrel{~}{C}_1`$ except finitely many points. Since $`𝒫_\omega `$ is superharmonic and $`𝒫_\omega (\xi )=𝒫_\mu (\xi )`$ for $`\xi \stackrel{~}{C}_1`$,

$$𝒫_\omega (z)\frac{1}{2\pi ϵ}_{|z\xi |=ϵ}𝒫_\omega (\xi )|d\xi |=\frac{1}{2\pi ϵ}_{|z\xi |=ϵ}𝒫_\mu (\xi )|d\xi |.$$

By letting $`ϵ0`$, we get $`𝒫_\omega (z)𝒫_\mu (z)`$ since $`𝒫_\mu (\xi )`$ is continuous at $`z\stackrel{~}{C}_p`$, $`p2`$. This completes the proof. ∎

###### Lemma 21.

Suppose $`(\stackrel{~}{C}_0)^{}\mathrm{}`$ and $`U`$ is a component of $`(\stackrel{~}{C}_0)^{}`$. Then there exists a subsequence of $`\{\mu _k\}`$ that converges in the weak-star topology to a measure $`\omega `$ such that

$$𝒫_\omega (z)=𝒫_\mu (z)\text{for all }z\overline{U}.$$

###### Proof.

Pick a point $`z_0U`$. By (1.6) there is a subsequence $`\{\mu _{k_j}\}`$ such that

$$\underset{j\mathrm{}}{lim}𝒫_{\mu _{k_j}}(z_0)=\underset{j\mathrm{}}{lim}k_j^1\mathrm{log}|P_{k_j}(z_0)|=\mathrm{log}\delta _0.$$

Now let $`\omega `$ be a weak-star limit of a subsequence of $`\{\mu _{k_j}\}`$. Then by the Principle of Descent (, Theorem 1.3), we have $`𝒫_\omega (z_0)\mathrm{log}\delta _0`$. Since (A.2) says that $`𝒫_\omega (z)=\mathrm{log}\delta _0`$ for all $`zU`$, the minimum principle implies that $`𝒫_\omega =\mathrm{log}\delta _0=𝒫_\mu `$ in $`U`$. ∎

###### Proofs of Theorems 5 and 6.

Suppose $`(\stackrel{~}{C}_0)^{}`$ is empty and let $`\omega `$ be a weak-star limit of any convergent subsequence of $`\{\mu _k\}`$. Then by (A.2), we have $`𝒫_\omega (z)=𝒫_\mu (z)`$ for all $`z`$. Therefore $`\omega =(2\pi )^1\mathrm{\Delta }𝒫_\omega =(2\pi )^1\mathrm{\Delta }𝒫_\mu =\mu `$ and the first statement of Theorem 5 follows from Lemma 19. If, in addition, $`=\stackrel{~}{C}_0\stackrel{~}{C}_1`$, then $`\mu `$ is the equilibrium distribution of $`\stackrel{~}{C}_0`$ since the support of $`\mu `$ is $`\stackrel{~}{C}_1=\stackrel{~}{C}_0`$ and $`𝒫_\mu `$ is constant on $`\stackrel{~}{C}_0`$ (, Theorem III.15). This completes the proof of Theorem 5.

Now suppose $`(\stackrel{~}{C}_0)^{}`$ is connected. Then by Lemma 21, there exists a subsequence $`\{\mu _{k_j}\}`$ of $`\{\mu _k\}`$ converging to a measure $`\omega `$ such that $`𝒫_\omega =𝒫_\mu `$ in $`\stackrel{~}{C}_0`$. Since $`𝒫_\omega (z)=𝒫_\mu (z)`$ for all $`z\backslash \stackrel{~}{C}_0`$ by (A.2), this shows that $`𝒫_\omega =𝒫_\mu `$ in $``$ hence $`\omega =\mu `$. Now Theorem 6 follows from Lemma 19. ∎


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# Commutator automorphisms of formal power series rings

## 1. Introduction

### 1.1. Motivation

A $`K`$-theoretical framework for some well known conjectures on automorphisms and idempotent endomorphisms of polynomial rings was initiated at the end of the 1970s in the works \[BW, C1, C2, CW\]. However, the development since then has been blocked due to intractability of these conjectures. A recent evidence of this intractability is provided by \[SU\].

There are two natural variations of the mentioned nonlinear $`K`$-theory that makes things easier and tractable. One variant corresponds to extending polynomial algebras to *polytopal algebras* – essentially affine cones over projective toric varieties – and restricting homomorphisms to graded homomorphisms. A systematic theory in this direction, including higher groups, is developed in \[BrG1, BrG2, BrG3, BrG4\].

Here we consider the other variation that corresponds to the completion process and leads to the category of formal power series and their continuous homomorphisms. For higher groups this approach leads to a challenging conjectural homological computation – the isomorphisms (1) below for which supporting low dimensional results are obtained in this paper.

The freeness of projective objects in the context of complete algebras was proved in \[T\], as an application of a technique designed to study retracts of polynomial algebras, and also in \[P\], in the general context of projective algebras over nilpotent algebraic theories. One should also mention the work \[K\], establishing an analogue of the $`K`$-theoretical local-global principle for complete algebras.

Below we show (Theorem 1.1) how the standard method of ‘indeterminate coefficients’ implies the aforementioned result on retracts of $`R[[X_1,\mathrm{},X_n]]`$. Furthermore, the same method almost immediately yields the corresponding $`K_1`$-analogue for the group $`\mathrm{Aut}(R[[X_1,\mathrm{},X_n]])`$ if the ground ring $`R`$ contains $``$ (Theorem 1.2).

Main emphasize in this paper is put on derivation of a similar result on the automorphism group when the coefficient ring is an algebra over a finite field (Theorem 1.3). Here the things are a bit complicated. One needs to invoke a complex induction on coefficient indices, depending on the residues modulo the characteristic. The main difficulty to make the indeterminate coefficients method work has been guessing many auxiliary automorphisms yielding feasible systems of equation. Also, one needs to treat separately the univariate case and the passage to the multivariate case. For a variant of the result in the case $`n=1`$ and $`R=_p`$, where $`p`$ is a prime number $`2,3`$, see \[Ca1\]. The corresponding groups of automorphisms are related to *Nottingham groups* – objects of recent active investigation \[Ca2\].

We have not been able to involve algebras over $`_2`$ and $`_3`$ in our consideration.

The higher $`K`$-theory machinery of symmetric monoidal categories \[Gr\] explains what the direct higher analogue of the main result (after interpreting it as a claim on the fundamental group of certain $`K`$-theoretic spaces) should be. We conjecture the isomorphisms of integral homologies

(1)
$$H_i(E(R),)=H_i([GA^c(R),GA^c(R)],),i2$$

where $`E(R)`$ is the stable group of elementary matrices and $`GA^c(R)`$ stands for the formal non-linear analogue of the general linear group $`GL(R)`$. These isomorphisms seem to be highly nontrivial already for $`i=2`$. In this particular case a theory of universal Steinberg relations should be possible, generalizing Milnor’s $`K_2`$. In the polytopal $`K`$-theory this is done in \[BrG3\].

### 1.2. Main results

Throughout the paper $`R`$ is a commutative ring.

$`\mathrm{Mod}(R)`$ denotes the category of $`R`$-modules, and $`(R)`$ denotes the full subcategory of finitely generated projective modules. Let $`\mathrm{Alg}^\mathrm{c}(R)`$ be the category of augmented $`R`$-algebras, complete in the additive topology of the augmentation ideal. For $`A\mathrm{Alg}^\mathrm{c}(R)`$ let $`I(A)=\mathrm{Ker}(AR)`$. The coproduct of two algebras in $`\mathrm{Alg}^\mathrm{c}(R)`$ is given by the additive completion of the tensor product: $`AB=(A_RB)_{I(AB)}^{}`$, $`I(AB)=\mathrm{Im}\left((I(A)B)(AI(B))AB\right)`$. This makes $`\mathrm{Alg}^\mathrm{c}(A)`$ a symmetric monoidal category. We have the functors:

$`\widehat{S}:\mathrm{Mod}(R)\mathrm{Alg}^\mathrm{c}(R),\widehat{S}(M)`$ $`=(RMS^2(M)\mathrm{})_{0MS^2(M)\mathrm{}}^{}`$
$`\mathrm{I}/\mathrm{I}^2:\mathrm{Alg}^\mathrm{c}(R)\mathrm{Mod}(R),\mathrm{I}/\mathrm{I}^2(A)`$ $`=I(A)/I^2(A).`$
Let $`\overline{\mathrm{X}}=(X_1,\mathrm{},X_n)`$. A *retract* of $`R[[\overline{\mathrm{X}}]]`$ is an algebra $`A\mathrm{Alg}^\mathrm{c}(R)`$ for which there are morphisms $`AR[[\overline{\mathrm{X}}]]A`$ whose composite is $`\mathrm{𝟏}_A`$.

###### Theorem 1.1.

Any retract of $`R[[\overline{\mathrm{X}}]]`$ is of type $`\widehat{S}(P)`$ for some $`P(R)`$.

Denote by $`GA_n^c(R)`$ the group of augmented continuous $`R`$-automorphisms of the formal power series ring $`R[[\overline{\mathrm{X}}]]`$, where ‘augmented’ and ‘continuous’ are understood with respect to the ideal $`(\overline{\mathrm{X}})R[[\overline{\mathrm{X}}]]R[[\overline{\mathrm{X}}]]`$ and the corresponding additive topology on $`R[[\overline{\mathrm{X}}]]`$. Such an automorphism is determined by its values on the variables

$$X_ih_i^{(1)}+h_i^{(2)}+\mathrm{},i=1,\mathrm{},n$$

where $`h_i^{(j)}R[\overline{\mathrm{X}}]`$ is a homogeneous polynomial of degree $`j`$ for $`i=1,\mathrm{},n`$.

The surjective group homomorphism $`GA_n^c(R)\mathrm{Aut}(R[[\overline{\mathrm{X}}]]/𝔐^2)=GL_n(R)`$ is split by the natural embedding $`GL_n(R)GA_n^c(R)`$ whose image consists of the linear transformations of the $`X_i`$. Let $`GI_n^c(R)`$ and $`GE_n^c(R)`$ denote respectively the preimages of the trivial and elementary subgroups: $`\{\mathrm{𝟏}\}E_n(R)GL_n(R)`$.

For a natural number $`k`$ and a group $`\mathrm{\Gamma }`$ we let $`[\mathrm{\Gamma },\mathrm{\Gamma }]^{(k)}`$ denote the subset

$$\{[\beta _1,\gamma _1]\mathrm{}[\beta _k,\gamma _k]:\beta _1,\gamma _1,\mathrm{},\beta _k,\gamma _k\mathrm{\Gamma }\}[\mathrm{\Gamma },\mathrm{\Gamma }]$$

where $`[,]`$ refers to the commutator.

###### Theorem 1.2.

If there exist $`b,c\mathrm{U}(R)`$ such that $`(b^m1)R+(c^m1)R=R`$ for all natural numbers $`m`$ then $`GI_n^c(R)[GA_n^c(R),GA_n^c(R)]^{(1)}`$.

The class of rings in Theorem 1.2 contains all $``$-algebras.

Let $`p`$ be a prime number $`5`$.

###### Theorem 1.3.

(a) If $`_pR`$ then $`GI_1^c(R)=[GA_1^c(R),GA_1^c(R)]^{(2)}`$.

(b) If $`_pR`$ and $`n2`$ then $`GI_n^c(R)[GA_n^c(R),GA_n^c(R)]^{(n(n+3))}`$.

(c) If $`R`$ is a field of characteristic $`p`$ then $`GI_n^c(R)[GA_n^c(R),GA_n^c(R)]^{(4n)}`$.

*Convention*. Below $`R^n`$ is thought of as the module of $`n`$-columns, $`\mathrm{U}(R)`$ refers to the group of units of $`R`$. For a univariate formal power series $`f(X)R[[X]]`$ its $`m`$-th coefficient is denoted by $`f(X)_m`$.

## 2. Proof of Theorem 1.1

Let $`AR[[\overline{\mathrm{X}}]]`$ be an $`R`$-subalgebra for which there exists an $`R`$-algebra homomorphism $`\pi :R[[\overline{\mathrm{X}}]]A`$ such that $`\pi |_A=\mathrm{𝟏}_A`$. We put

$$\pi (X_i)=f_i=f_i^{(1)}+f_i^{(2)}+\mathrm{},i=1,\mathrm{},n$$

where $`f_i^{(k)}R[\overline{\mathrm{X}}]`$ is a homogeneous polynomial of degree $`k`$. Let $`M(\pi )`$ denote the $`n\times n`$-matrix whose $`i`$th column is $`(a_{i1},\mathrm{},a_{in})^{}`$ where $`f_i^{(1)}=a_{i1}X_1+\mathrm{}+a_{in}X_n`$. The submodule of $`R^n`$ generated by the columns of $`M(\pi )`$ is naturally identified with the module $`P=(I/I^2)(A)`$. The latter is a projective $`R`$-module: we have the split epimorphism $`(I/I^2)(\pi ):R^nP`$. Put $`Q=\mathrm{Ker}\left((I/I^2)(\pi )\right)`$ and fix an epimorphism $`\rho :R^nQ`$ split by the embedding $`QR^n`$. Then we have the following split $`R`$-algebra epimorphism

$$\pi \widehat{S}(\rho ):R[[\overline{\mathrm{X}},\overline{\mathrm{Y}}]]A\widehat{S}(Q).$$

(Here $`\overline{\mathrm{Y}}=(Y_1,\mathrm{},Y_n)`$.) Because $`PQR^n`$ there exists $`\alpha \mathrm{Aut}_R(R^{2n})=GL_{2n}(R)`$ such that $`\alpha (PQ)=R^n0`$. Then the idempotent endomorphism $`\kappa =\widehat{S}(\alpha )\left(\pi \widehat{S}(\rho )\right)\widehat{S}(\alpha )^1:R[[\overline{\mathrm{X}},\overline{\mathrm{Y}}]]R[[\overline{\mathrm{X}},\overline{\mathrm{Y}}]]`$ is of the form

$`\kappa (X_i)=X_i+g_i^{(2)}+g_i^{(3)}+\mathrm{},i=1,\mathrm{},n,`$
$`\kappa (Y_j)=h_j^{(2)}+h_j^{(3)}+\mathrm{},j=1,\mathrm{},n`$
where $`g_i^{(k)},h_j^{(k)}R[\overline{\mathrm{X}},\overline{\mathrm{Y}}]`$ are homogeneous polynomials of degree $`k`$. We want to show the equality

(2)
$$\mathrm{Im}(\kappa )=R[[\overline{\mathrm{X}}]].$$

First observe that $`h_j^{(2)}R[\overline{\mathrm{X}}]`$, $`j=1,\mathrm{},n`$. In fact, if $`h_j^{(2)}R[\overline{\mathrm{X}},\overline{\mathrm{Y}}]R[\overline{\mathrm{X}}]`$ for some index $`j`$ then the first appearance of the $`Y`$-variables in $`\kappa ^2(Y_j)`$ can only be in a homogeneous summand of degree $`>2`$, contradicting the condition $`\kappa ^2=\kappa `$.

Assume we have shown

(3)
$$h_j^{(2)},h_j^{(3)},\mathrm{},h_j^{(k)}R[\overline{\mathrm{X}}],j=1,\mathrm{},n.$$

Applying $`\kappa ^2=\kappa `$ to the variables $`X_i`$ we get

(4)
$$g_i^{(2)},g_i^{(3)},\mathrm{},g_i^{(k)}R[\overline{\mathrm{X}}],i=1,\mathrm{},n.$$

Therefore, (2) follows once it is shown that (3) implies

$$h_j^{(k+1)}R[\overline{\mathrm{X}}],j=1,\mathrm{},n.$$

Assume to the contrary that $`h_j^{(k+1)}`$ involves one of the variables $`Y_1,\mathrm{},Y_n`$ for some index $`j`$. But then, in view of (3) and (4), the first appearance of the $`Y`$-variables in $`\kappa ^2(Y_j)`$ can only be in a homogeneous summand of degree higher than $`k+1`$, contradicting the condition $`\kappa ^2=\kappa `$.

The following commutative diagram is a consequence of (2):

where the vertical arrows represent isomorphisms, the composite of the first two upper homomorphisms is $`\pi |_{\widehat{S}(P)}\mathrm{𝟏}_{\widehat{S}(Q)}`$, and that of the first two lower homomorphisms is an isomorphism. It follows that $`\pi |_{\widehat{S}(P)}:\widehat{S}(P)A`$ is also an isomorphism.∎

## 3. Composite automorphisms

Here we prove Theorem 1.2. Let $`n`$ be a fixed natural number.

Assume $`\beta ,\gamma GA_1^c(R)=\mathrm{Aut}(R[[X]])`$, $`\beta (X)=b_1X+b_2X^2+\mathrm{}`$ and $`\gamma (X)=c_1X+c_2X^2+\mathrm{}`$. We have

(5)
$$(\gamma \beta )(X)_m=\underset{r=1}{\overset{m}{}}b_r\underset{\begin{array}{c}(j_1,\mathrm{},j_r)\\ j_1+\mathrm{}+j_r=m\end{array}}{}c_{j_1}\mathrm{}c_{j_r}.$$

While there is no obvious compact multivariate analog of (5) there is some sort of control on composite automorphisms which will play crucial rôle in the sequel.

For a natural number $`m`$ let $`𝒮_m`$ denote the number of ordered partitions of $`m`$ with nonnegative entries and having length $`n`$. Arbitrary continuous $`R`$-algebra endomorphism $`\alpha `$ of $`R[[\overline{\mathrm{X}}]]`$ can be represented in the following way:

$`\alpha (\overline{\mathrm{X}})={\displaystyle \underset{m=1}{\overset{\mathrm{}}{}}}\alpha (\overline{\mathrm{X}})_m\left((X_1^{d_1}\mathrm{}X_n^{d_n})_{d_1+\mathrm{}+d_n=m}\right)^{}`$
where $`\alpha (\overline{\mathrm{X}})_m`$ is an $`n\times 𝒮_m`$ matrix over $`R`$ and

$$\left((X_1^{d_1}\mathrm{}X_n^{d_n})_{d_1+\mathrm{}+d_n=m}\right)$$

refers to the vector of monomials of total degree $`d`$ whose components are ordered lexicographically with respect to $`X_1>\mathrm{}>X_n`$.

Next we introduce the pairings

$$M_{n\times 𝒮_l}(R)\times M_{n\times 𝒮_m}(R)M_{n\times 𝒮_{lm}}(R),(L,M)LM,l,m,$$

defined by

$`\left(L\left((X_1^{d_1}\mathrm{}X_n^{d_n})_{d_1+\mathrm{}+d_n=l}\right)^{}\right)\left(M\left((X_1^{d_1}\mathrm{}X_n^{d_n})_{d_1+\mathrm{}+d_n=m}\right)^{}\right)\left(\overline{X}\right)=`$
$`(LM)\left((X_1^{d_1}\mathrm{}X_n^{d_n})_{d_1+\mathrm{}+d_n=lm}\right)`$
where $``$ is the composition of the two endomorphisms of $`R[[\overline{X}]]`$. The associativity of the composition yields the associativity rule $`(LM)N=L(MN)`$.

Clearly, when $`l=m=1`$ the $``$ operation becomes the usual matrix product of $`n\times n`$ matrices.

For natural numbers $`r`$ and $`m`$ and a system of automorphisms $`\delta _1,\mathrm{},\delta _rGA_n^c(R)`$ the components of the matrix $`(\delta _1\mathrm{}\delta _k)(\overline{\mathrm{X}})_m`$ are polynomial functions of the components of the matrices $`\delta _1(\overline{\mathrm{X}})_s,\mathrm{},\delta _k(\overline{\mathrm{X}})_s`$, $`s=1,2,\mathrm{}`$, the latter being treated as variables in the next lemma.

###### Lemma 3.1.

Let $`\alpha ,\beta ,\gamma `$ be continuous $`R`$-algebra endomorphisms of $`R[[\overline{X}]]`$.

* The components of the matrices $`(\beta \gamma \alpha )(\overline{\mathrm{X}})_k`$ and $`(\gamma \beta )(\overline{\mathrm{X}})_k`$ do not depend on the components of the matrices $`\alpha (\overline{\mathrm{X}})_m`$, $`\beta (\overline{\mathrm{X}})_m`$, or $`\gamma (\overline{\mathrm{X}})_m`$ for $`k<m`$.
* The components of the $`n\times 𝒮_m`$ matrices

$$(\beta \gamma \alpha )(\overline{\mathrm{X}})_m\beta (\overline{\mathrm{X}})_m\gamma (\overline{\mathrm{X}})_1\alpha (\overline{\mathrm{X}})_1\beta (\overline{\mathrm{X}})_1\gamma (\overline{\mathrm{X}})_m\alpha (\overline{\mathrm{X}})_1$$

and

$$(\gamma \beta )(\overline{\mathrm{X}})_m\gamma (\overline{\mathrm{X}})_m\beta (\overline{\mathrm{X}})_1\gamma (\overline{\mathrm{X}})_1\beta (\overline{\mathrm{X}})_m$$

do not depend on the components of the matrices $`\beta (\overline{\mathrm{X}})_m`$ or $`\gamma (\overline{\mathrm{X}})_m`$.

Proof is straightforward.

Below for an element $`\alpha GI_n^c(R)`$ we will consider the equation $`\alpha =[\gamma ^1,\beta ^1]`$ to be solved for $`\beta ,\gamma GA_n^c(R)`$. This is equivalent to the infinite system of equations

($`_m`$) $`(\beta \gamma \alpha )(\overline{\mathrm{X}})_m=(\gamma \beta )(\overline{\mathrm{X}})_m,m,`$
for matrices $`\beta (\overline{\mathrm{X}})_m,\gamma (\overline{\mathrm{X}})_mM_{n\times 𝒮_m}(R)`$ such that $`\beta (\overline{\mathrm{X}})_1,\gamma (\overline{\mathrm{X}})_1GL_n(R)`$.

###### Proof of Theorem 1.2.

Let $`b,cR`$ be elements such that $`(b^m1)R+(c^m1)R=R`$ for all natural numbers $`m`$. For arbitrary element $`\alpha GI_n^c(R)`$ we want to solve the infinite system $`(_m)`$, $`m`$.

Let $`\beta (\overline{\mathrm{X}})_1=b\mathrm{𝐈𝐝}_n`$ and $`\gamma (\overline{\mathrm{X}})_1=c\mathrm{𝐈𝐝}_n`$. It follows from Lemma 3.1(a) that whatever matrices $`\beta (\overline{\mathrm{X}})_m`$ and $`\gamma (\overline{\mathrm{X}})_m`$, $`m2`$ we take (of size $`n\times 𝒮_m`$) the equality $`(_1)`$ is satisfied. More generally, the same proposition implies that for arbitrary natural number $`m2`$ the validity of $`(_k)`$, $`k<m`$ only depends on the matrices $`\beta (\overline{\mathrm{X}})_k`$ and $`\gamma (\overline{\mathrm{X}})_k`$, $`k<m`$.

Now Lemma 3.1(b) implies that for every natural number $`m2`$ the equation $`(_m)`$ rewrites as

$`(\beta \gamma \alpha )(\overline{\mathrm{X}})_m(\gamma \beta )(\overline{\mathrm{X}})_m=`$ $`\beta (\overline{\mathrm{X}})_1\gamma (\overline{\mathrm{X}})_m\alpha (\overline{\mathrm{X}})_1+\beta (\overline{\mathrm{X}})_m\gamma (\overline{\mathrm{X}})_1\alpha (\overline{\mathrm{X}})_1`$
$`\gamma (\overline{\mathrm{X}})_m\beta (\overline{\mathrm{X}})_1\gamma (\overline{\mathrm{X}})_1\beta (\overline{\mathrm{X}})_m+_m`$
where $`_m`$ is an $`n\times 𝒮_m`$ matrix which only depends on the entries of the matrices $`\alpha (\overline{\mathrm{X}})_k,\beta (\overline{\mathrm{X}})_k,\gamma (\overline{\mathrm{X}})_k`$, $`k<m`$. It is, therefore, possible to find successively matrices $`\beta (\overline{\mathrm{X}})_m`$ and $`\gamma (\overline{\mathrm{X}})_m`$ with the desired properties once the following is shown – for arbitrary natural number $`m`$ and arbitrary $`n\times 𝒮_m`$ matrix $`𝒜`$ there are two $`n\times 𝒮_m`$ matrices $``$ and $`𝒞`$ such that

$$𝒜=\beta (\overline{\mathrm{X}})_1𝒞\alpha (\overline{\mathrm{X}})_1+\gamma (\overline{\mathrm{X}})_1\alpha (\overline{\mathrm{X}})_1𝒞\beta (\overline{\mathrm{X}})_1\gamma (\overline{\mathrm{X}})_1.$$

In view of the conditions $`\alpha (\overline{\mathrm{X}})_1=\mathrm{𝐈𝐝}`$, $`\beta (\overline{\mathrm{X}})_1=b\mathrm{𝐈𝐝}`$ and $`\gamma (\overline{\mathrm{X}})_1=c\mathrm{𝐈𝐝}`$ this equation is equivalent to $`𝒜=𝒞(b^mb)+(cc^m)`$. Finally, the existence of the desired matrices $``$ and $`𝒞`$ follows from the condition

$$(b^{m1}1)(bR)+(c^{m1}1)(cR)=(b^{m1}1)R+(c^{m1}1)R=R$$

being applied separately to every component. ∎

The proof of Theorem 1.2 implies the following useful fact to be used later on:

###### Corollary 3.2.

Let $`m`$ be a natural number. If there exist elements $`b,c\mathrm{U}(R)`$ such that $`(b^k1)R+(c^k1)R=R`$ for all $`k=1,\mathrm{},m1`$ then for arbitrary element $`\alpha GI_n^c(R)`$ there exist elements $`\beta ,\gamma GA_n^c(R)`$ such that $`\alpha (\overline{\mathrm{X}})_k=[\beta ,\gamma ](\overline{\mathrm{X}})_k`$ for $`k=1,\mathrm{},m`$.

## 4. characteristic $`5`$: the univariate case

In this section we prove Theorem 1.3(a). The notation $`(_m)`$ has the same meaning as in Section 3, only considered in the univariate case.

Let $`𝐅`$ be a field contained in $`R`$, $`\mathrm{char}𝐅=p5`$.

Choose arbitrary elements $`b_2,b_3,c_2,c_3𝐅`$ such that $`c_2b_20`$. Put $`a_2=[\gamma ^1,\beta ^1](X)_2`$ and $`a_3=[\gamma ^1,\beta ^1](X)_3`$ where $`\beta ,\gamma GA_1^c(R)`$ are arbitrary automorphisms such that $`\beta (X)_1=\gamma (X)=1`$ and $`\beta (X)_k=b_k`$, $`\gamma (X)_k=c_k`$, $`k=2,3`$. Using (5) in Section 3 one easily sees $`a_2=2(b_2c_2)`$ and $`a_3=4(b_2c_2)^2`$. In particular, if $`b_2=1`$, $`c_2=0`$ then $`a_2=2`$ and $`a_3=4`$.

###### Lemma 4.1.

Let $`\alpha GA_1^c(R)`$ arbitrary element such that $`\alpha (X)=X+a_2X^2+a_3X^3+a_4X^4+\mathrm{}`$. Assume $`m4`$ is an even natural number, $`m2mod(p)`$. Then for arbitrary elements $`b_k,c_kR`$, $`k4`$, $`km`$ there exist elements $`b_m,c_mR`$ such that $`(_m)`$ and $`(_{m+1})`$ are satisfied for the automorphisms $`\beta ,\gamma GA_1^c(R)`$ where $`\beta (X)=X+b_2X^2+b_3X^3+b_4X^4+\mathrm{}`$ and $`\gamma (X)=X+c_2X^2+c_3X^3+c_4X^4+\mathrm{}`$.

###### Proof.

Using (5) and the conditions on $`m`$ one easily checks that the system $`(_m,_{m+1})`$ is independent of $`b_{m+1}`$ and $`c_{m+1}`$ (although the values of the $`(m+1)`$-coefficients do depend on $`b_{m+1}`$ and $`c_{m+1}`$). Moreover, this system rewrites as the following system of linear equations for $`c_m`$ and $`b_m`$:

(A) $`2`$ $`c_m`$ $`2`$ $`b_m`$ $`=𝒳`$
$`\left(+2a_2(m2)b_2\right)`$ $`c_m+`$ $`\left(2a_2+(m2)c_2\right)`$ $`b_m`$ $`=𝒴`$
where $`𝒳`$ and $`𝒴`$ are elements of $`R`$ only depending (polynomially) on $`b_k,c_k`$, $`k<m`$. See Remark 4.3 below. Now the determinant of $`(𝐀)`$ is $`2(m2)(c_2b_2)\mathrm{U}(R)`$ and, therefore, the system is (uniquely) solvable for $`b_m`$ and $`c_m`$. ∎

###### Lemma 4.2.

Assume $`m4`$ is an even natural number, $`m=2mod(p)`$. Let $`b_2=1`$, $`c_2=0`$, and $`b_3,c_3𝐅`$ such that $`b_3+c_3=0`$. Then for arbitrary elements $`b_k,c_kR`$, $`k4`$, $`km1,m`$ there exist elements $`b_{m1},c_{m1},b_m,c_mR`$ such that $`(_m)`$ and $`(_{m+1})`$ are satisfied for the automorphisms $`\beta ,\gamma GA_1^c(R)`$ where $`\beta (X)=X+b_2X^2+b_3X^3+b_4X^4+\mathrm{}`$ and $`\gamma (X)=X+c_2X^2+c_3X^3+c_4X^4+\mathrm{}`$.

###### Proof.

We start with general elements $`b_2,c_2,b_3,c_3𝐅`$. By multiple applications of (5) the system $`(_m,_{m+1})`$ rewrites as

($`𝐁`$) $`2(c_mb_m)+\left(2a_2+(m+1)b_2\right)c_{m1}+\left(2a_2(m+1)c_2\right)b_{m1}`$ $`=𝒳`$
$`2a_2(c_mb_m)+\left(3a_3+(m4)b_3+2ma_2b_22a_2c_22b_2c_2\right)c_{m1}`$ $`+`$
$`+\left(3a_3(m4)c_3+2a_2b_26a_2c_2+2c_2b_2\right)b_{m1}`$ $`=𝒴`$
where $`𝒳,𝒴R`$ only depend (polynomially) on $`b_4,c_4,\mathrm{},b_{m2},c_{m2}`$. See Remark 4.3 below. We view $`(𝐁)`$ as a system of linear equations with respect to $`2(c_mb_m)`$, $`c_{m1}`$, and $`b_{m1}`$. Since the matrix of $`(𝐁)`$ is defined over $`𝐅`$ it is enough to show that its rank equals 2. By multiplying the first row of the mentioned matrix by $`a_2`$ and then subtracting the result from the second row the condition on the rank becomes $`(m3)a_2(b_2c_2)4b_2c_2+(m4)(b_3+c_3)0`$. For our specific choice of $`b_2,c_2,b_3,c_3`$ this is equivalent to $`m3mod(p)`$, which is the case as $`m=2mod(p)`$. ∎

###### Remark 4.3.

Here we explain how the coefficients in the systems $`(𝐀)`$ and $`(𝐁)`$ above are computed. We only illustrate the process on the coefficients of $`c_{m1}`$ and $`b_{m1}`$ in the second equation of $`(𝐁)`$. The argument for all other coefficients is similar and substantially easier.

Put $`(\beta \gamma )(X)=X+d_2X^2+d_3X^3+\mathrm{}`$, the $`d_j`$ being polynomial functions of $`b_2,c_2,b_3,c_3,\mathrm{}`$. Then by (5) we have

$$(\beta \gamma \alpha )(X)_{m+1}=\underset{r=1}{\overset{m+1}{}}a_r\underset{\begin{array}{c}(j_1,\mathrm{},j_r)\\ j_1+\mathrm{}+j_r=m+1\end{array}}{}d_{j_1}\mathrm{}d_{j_r},$$

where $`d_1=1`$. By Lemma 3.1(a) among the $`d_j`$ in this equality only $`d_{m1}`$, $`d_m`$ and $`d_{m+1}`$ depend on the coefficients $`c_{m1}`$ and $`b_{m1}`$. Thus, $`(\beta \gamma \alpha )(X)_{m+1}=(2d_2a_2+3a_3)d_{m1}+2a_2d_m+d_{m+1}+𝒰`$ where $`𝒰`$ does not depend on $`c_{m1},b_{m1}`$.

Again (5), together with the assumptions that $`m`$ is even and $`b_1=c_1=1`$, implies $`d_2=c_2b_2`$, $`d_{m1}=(c_{m1}b_{m1})+𝒰_{m1}`$, $`d_m=(m1)b_2c_{m1}2c_2b_{m1}+𝒰_m`$ and $`d_{m+1}=\left((m1)b_3\left(\genfrac{}{}{0pt}{}{m1}{2}\right)b_2^2\right)c_{m1}+\left(2b_2c_2+3c_3\right)b_{m1}+𝒰_{m+1}`$ where $`𝒰_{m1}`$, $`𝒰_m`$ and $`𝒰_{m+1}`$ do not depend on $`c_{m1}`$ and $`b_{m1}`$. Also, the condition $`m=2mod(p)`$ implies $`\left(\genfrac{}{}{0pt}{}{m1}{2}\right)=0R`$.

Similar computations show $`(\gamma \beta )(X)_{m+1}=\left(3b_3+2b_2c_2\right)c_{m1}+(m1)c_3b_{m1}+𝒲`$ where $`𝒲`$ does not depend on $`c_{m1}`$ and $`b_{m1}`$. Summing up, we get the desired coefficients in the expression

$`(\beta \gamma \alpha )(X)_{m+1}`$ $`(\gamma \beta )(X)_{m+1}=`$
$`\left(3a_3+(m4)b_3+2ma_2b_22a_2c_22b_2c_2\right)c_{m1}+`$
$`\left(3a_3(m4)c_3+2a_2b_26a_2c_2+2c_2b_2\right)b_{m1}+𝒱+𝒲`$
###### Proof of Theorem 1.3(a).

Let $`\alpha GI_1^c(R)`$ be an arbitrary automorphism and $`\phi GA_1^c(R)`$ be the automorphism $`\phi (X)=X+2X^2+4X^3`$. Since $`\mathrm{char}𝐅5`$ the field $`𝐅`$ contains an element of order $`4`$. Then by Corollary 3.2 (and Lemma (3.1)) there exist elements $`\beta _0,\gamma _0GA_1^c(R)`$ such that $`\left(\alpha \phi ^1\right)(X)_k=[\beta _0,\gamma _0](X)_k`$ for $`k=2,3`$, equivalently $`\left([\beta _0,\gamma _0]^1\alpha \right)(X)_2=2`$ and $`\left([\beta _0,\gamma _0]^1\alpha \right)(X)_3=4`$. We are done because by Lemmas 4.1 and 4.2 any automorphism $`\psi GI_1^c(R)`$, such that $`\psi (X)_2=2`$ and $`\psi (X)_3=4`$, belongs to the commutator subgroup $`[GA_1^c(R),GA_1^c(R)]^{(1)}`$. ∎

## 5. characteristic $`5`$: the multivariate case

In this section we assume $`n2`$.

For a power series $`g(X_1,\mathrm{},\stackrel{ˇ}{X}_i,\mathrm{},X_n)^2R[[X_1,\mathrm{},\stackrel{ˇ}{X}_i,\mathrm{},X_n]]`$ the element $`\epsilon _i^g\mathrm{Aut}(R[\overline{\mathrm{X}}])`$, defined by

$$\epsilon _i^g(X_j)=\{\begin{array}{cc}X_i+g\text{for}j=i,\hfill & \\ X_j\text{for}ji,\hfill & \end{array}$$

will be called *elementary*.

###### Lemma 5.1.

Elementary automorphisms of $`R[[\overline{\mathrm{X}}]]`$ belong to $`[GA_n^c(R),GA_n^c(R)]^{(1)}`$.

###### Proof.

It is enough to observe that $`\epsilon _1^g[GA_n^c(R),GA_n^c(R)]^{(1)}`$ for arbitrary element $`g(X_2,\mathrm{},X_n)^2R[[X_2,\mathrm{},X_n]]`$. In fact, we have $`\beta \epsilon _1^g\gamma =\gamma \beta `$ where

$$\beta (X_i)=\{\begin{array}{cc}X_1+X_2\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n\hfill & \end{array}\text{and}\gamma (X_i)=\{\begin{array}{cc}X_i\text{for}i=1,3,\mathrm{},n,\hfill & \\ X_2+g\text{for}i=2.\hfill & \end{array}$$

###### Lemma 5.2.

Let $`\frac{1}{2}R`$ and $`f(X_2,\mathrm{},X_n)R[[X_2,\mathrm{},X_n]]`$. Then the automorphism $`\alpha GE_n^c(R)`$, defined by

$$\alpha (X_i)=\{\begin{array}{cc}X_1+X_1f\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n,\hfill & \end{array}$$

belongs to $`[GA_n^c(R),GA_n^c(R)]^{(n)}`$. Moreover, if $`R`$ is a field of characteristic $`2`$ then $`\alpha [GA_n^c(R),GA_n^c(R)]^{(1)}`$.

###### Proof.

*Step 1.* First consider the case when $`f(X_2,\mathrm{},X_n)^2R[[X_2,\mathrm{},X_n]]`$. In this situation we have $`\alpha \beta \gamma =\gamma \beta `$ where the elements $`\beta ,\gamma GA_n^c(R)`$ are defined by

$$\beta (X_i)=\{\begin{array}{cc}X_1+X_1X_2\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n\hfill & \end{array}\text{and}\gamma (X_i)=\{\begin{array}{cc}X_i\text{for}i=1,3,\mathrm{},n,\hfill & \\ X_2+f+X_2f\text{for}i=2.\hfill & \end{array}$$

In particular, $`\alpha [GA_n^c(R),GA_n^c(R)]^{(1)}`$.

*Step 2.* Now assume $`f=\xi X_2`$ for some $`\xi R`$. Then $`\alpha ^2\gamma =\gamma \alpha `$ where the automorphism $`\gamma GA_n^c(R)`$ is given by

$$\gamma (X_i)=\{\begin{array}{cc}X_i\text{for}i=1,3,4,\mathrm{},n,\hfill & \\ 2X_2+\xi X_2^2fori=2.\hfill & \end{array}$$

In particular, $`\alpha [GA_n^c(R),GA_n^c(R)]^{(1)}`$.

*Step 3.* In the general case $`\xi _2X_2+\mathrm{}+\xi _nX_n+g`$ for some $`\xi _2,\mathrm{},\xi _nR`$ and $`g(X_2,\mathrm{},X_n)^2R[[X_2,\mathrm{},X_n]]`$. By iterated use of the previous step (w.r.t. the variables $`X_2,\mathrm{},X_n`$) the automorphism $`\alpha ^{}GA_n^c(R)`$, defined by

$$\alpha ^{}(X_i)=\{\begin{array}{cc}X_1(1+\xi _2X_2)(1+\xi _3X_3)\mathrm{}(1+\xi _nX_n)\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n,\hfill & \end{array}$$

belongs to $`[GA_n^c(R),GA_n^c(R)]^{(n1)}`$. Finally, $`\alpha =\alpha ^{\prime \prime }\alpha ^{}`$ for the element $`\alpha ^{\prime \prime }GA_n^c(R)`$, defined by

$$\alpha ^{\prime \prime }(X_i)=\{\begin{array}{cc}X_1\vartheta \text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n,\hfill & \end{array}$$

where

$$\vartheta =\frac{1+\xi _2X_2+\mathrm{}+\xi _nX_n}{(1+\xi _2X_2)(1+\xi _3X_3)\mathrm{}(1+\xi _nX_n)}1+(X_2,\mathrm{},X_n)^2R[[X_2,\mathrm{},X_n]].$$

Now the case of rings containing 1/2 follows by Step 1.

*Step 4.* Consider the remaining case when $`R`$ is a field. By Step 1 there is no loss of generality in assuming $`f=\xi _2X_2+\mathrm{}+x_nX_n+g`$ where $`\xi _2\mathrm{U}(R)`$ and $`g(X_2,\mathrm{},X_n)^2R[[X_2,\mathrm{},X_n]]`$.

If $`\xi _21`$ then the exact same argument as in Step 1 works. Therefore, we can additionally assume that $`\xi _2=1`$. But then $`\alpha \beta _1\gamma _1=\gamma _1\beta _1`$ for the automorphisms $`\beta _1,\gamma _1GA_n^c(R)`$ determined by

$$\beta _1(X_i)=\{\begin{array}{cc}X_1X_1X_2\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n,\hfill & \end{array}\text{and}\gamma _1(X_i)=\{\begin{array}{cc}X_i\text{for}i=1,3,4,\mathrm{},n,\hfill & \\ X_2+(1+X_2)f\text{for}i=2.\hfill & \end{array}$$

###### Lemma 5.3.

Assume $`_pR`$ and $`h=(\overline{\mathrm{X}})^2R[[\overline{\mathrm{X}}]]`$. Then any automorphism $`\alpha GA_n^c(R)`$ of type

$$\alpha (X_1)=\{\begin{array}{cc}X_1+h\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n\hfill & \end{array}$$

belongs to $`[GA_n^c(R),GA_n^c(R)]^{(n+3)}`$. If $`R`$ is a field of characteristic $`p`$ then $`\alpha [GA_n^c(R),GA_n^c(R)]^{(4)}`$.

###### Proof.

We have $`h=g+X_1f_1+X_1^2f_2+X_1^3f_3+\mathrm{}`$ for some

$`g(X_2,\mathrm{},X_n)^2R[[X_2,\mathrm{},X_n]],f_1(X_2,\mathrm{},X_n)R[[X_2,\mathrm{},X_n]],`$
$`f_2,f_3,\mathrm{}R[[X_2,\mathrm{},X_n]].`$
First assume $`g=0`$. Then $`\alpha =\alpha ^{\prime \prime }\alpha ^{}`$ for the automorphisms $`\alpha ^{},\alpha ^{\prime \prime }GA_n^c(R)`$ defined by

$$\alpha ^{}(X_i)=\{\begin{array}{cc}X_1+X_1f_1\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n\hfill & \end{array},\alpha ^{\prime \prime }(X_i)=\{\begin{array}{cc}X_1+X_1^2\vartheta _2+\mathrm{}\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n\hfill & \end{array}$$

where $`\vartheta _j=(1+f_1)^1f_j`$, $`j=2,3,\mathrm{}`$. So Theorem 1.3(a) and Lemma 5.2 imply $`\alpha [GA_n^c(R),GA_n^c(R)]^{(n+2)}`$ in general and $`\alpha [GA_n^c(R),GA_n^c(R)]^{(3)}`$ for $`R`$ a field.

If $`g0`$ then $`\alpha =\beta \epsilon _1^g`$ where $`\beta GA_n^c(R)`$ is an automorphism of the type considered in the preceding paragraph. So by Lemma 5.1 we are reduced to the previous case $`g=0`$. ∎

###### Proof of Theorem 1.3(b,c).

Consider an automorphism

$$\alpha GI_n^c(R),\alpha (X_i)=X_i+h_i,h_i(\overline{\mathrm{X}})^2R[[\overline{\mathrm{X}}]],i=1,\mathrm{},n.$$

We have the equality $`\alpha =\alpha _1\alpha _2\mathrm{}\alpha _n`$ for the recursively defined elements $`\alpha _iGA_n^c(R)`$:

$`\alpha _1(X_i)`$ $`=\{\begin{array}{cc}X_1+h_1\text{for}i=1,\hfill & \\ X_i\text{for}i=2,\mathrm{},n,\hfill & \end{array}`$
$`\alpha _2(X_i)`$ $`=\{\begin{array}{cc}X_i\text{for}i=1,3,4,\mathrm{},n,\hfill & \\ \alpha _1^1\left(X_2+h_2\right)\text{for}i=2,\hfill & \end{array}`$
$`\alpha _3(X_i)`$ $`=\{\begin{array}{cc}X_i\text{for}i=1,2,4,5,\mathrm{},n,\hfill & \\ \alpha _2^1\alpha _1^1\left(X_3+h_3\right)\text{for}i=3,\hfill & \end{array}`$
$`\mathrm{}`$ $`\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$`\alpha _n(X_i)`$ $`=\{\begin{array}{cc}X_i\text{for}i=1,2,\mathrm{},n1,\hfill & \\ \alpha _{n1}^1\mathrm{}\alpha _2^1\alpha _1^1\left(X_n+h_n\right)\text{for}i=n.\hfill & \end{array}`$
It only remains to notice that each of these automorphisms is (up to enumeration of variables) of the type considered in Lemma 5.3. ∎


text:
# The distance to IC4499 from 𝐾-band photometry of 32 RR Lyrae starsBased on data from the European Southern Observatory

## 1 Introduction

The globular cluster IC4499 (RA: $`15^h00^m19^s`$, DEC: $`82^{}12^{}50^{\prime \prime }`$; 2000.0) is among the youngest galactic globular clusters (Ferraro et al. Ferraro95 (1995), Fusi Pecci et al., FusiPecci95 (1995)). It is of Oosterhoff type I and is thus likely to have zero or even retrograde rotation around the Galactic center (van den Bergh vdBergh93 (1993), Lee and Carney Lee99 (1999)). It is a fairly low density cluster with one of the highest numbers of RR Lyrae stars per unit luminosity (Suntzeff, Kinman and Kraft Suntzeff91 (1991)) containing about 100 RR Lyrae stars. These extreme characteristics suggests that the cluster might be a late addition to the Galaxy, which makes the cluster particularly interesting to study.

Optical color-magnitude diagrams have been presented by Sarajedini (Sarajedini93 (1993)) and Ferraro et al. (Ferraro95 (1995)), the latter extending two magnitudes below the main-sequence turn-off. The RR Lyrae stars have been the subject of several investigations (Coutts et al. Coutts75 (1975), Clements et al. Clement79 (1979), Clement86 (1986), Walker and Nemec WN96 (1996) (WN96 in the following)) since their discovery (Fourcade and Laborde, Fourcade69 (1969), and Fourcade et al., Fourcade74 (1974)). These investigations have revealed that the RR Lyrae stars span the full range of pulsation modes currently known, including stars exhibiting double mode pulsations, which can be used to determine the mass of these stars. WN96 has even suggested that some of the stars might pulsate in the second overtone. The range of the stellar masses has been shown by Sarajedini (Sarajedini93 (1993)) to be very limited as judged from the double mode pulsators.

The reddening towards IC4499 is fairly high, $`E(BV)=0.22\pm 0.02`$ (WN96) and thus a distance estimate which is insensitive to reddening will provide an important additional constraint on the distance to this cluster.

The $`M_K`$\- $`\mathrm{log}P`$ relation for RR Lyrae stars, pioneered by Longmore et al. (Longmore86 (1986)) and refined further by Longmore et al. (Longmore90 (1990)), is a very powerful tool for distance determination. The method is theoretically based on the pulsation equation from van Albada and Baker (AB71 (1971)) which links the period to the luminosity of an object. Bono et al. (Bono01 (2001)), and most recently Bono et al. (Bono03 (2003)) (B03 hereinafter) has refined the method using pulsational stellar models to include a small but significant metallicity term, and we will use this calibration in the following. It has been shown by Butler (Butler03 (2003)), through a comparison of inner and outer region RR Lyrae stars in M3, that the method is insensitive to environmental effects. It is also rather insensitive to errors in both \[Fe/H\] and reddening, and as the amplitudes of RR Lyrae stars in the $`K`$-band are smaller than about 0.3 mag it is feasible to apply the relation to single epoch data if a sufficiently large number of stars can be observed. Using recent ephemerides for the IC4499 stars we can even phase the observations and use template light-curves to improve our estimate of the mean $`K`$-magnitudes.

## 2 The Data

We used the IRAC-2B near-IR camera (Moorwood et al. Moorwood92 (1992)) mounted at the ESO-MPI 2.2m telescope at La Silla on the night Sep.17/18, 1995. The detector was a 256 by 256 pixel NICMOS-3 array which was read out using the so called ”non-destructive read” mode. The lens LC, which provides a projected pixel size of 0.5 arcsec was employed, and the seeing was typically 1.3 arcsec FWHM.

Two adjacent fields (labeled A and B in the following) including the center of IC4499 were selected on the basis of the charts in Clement et al. (Clement86 (1986)) to cover a large number of known RR Lyrae variables as well as a good range in period.

The observations were performed in a regular pattern of two dithered target frames interspersed by two dithered sky frames (several arc-minutes away) followed by the same sequence again, after offsetting the telescope by a few arc-seconds. Each frame was made of six 10 seconds integrations which were averaged on-line by the pre-processing system. For field A, 16 individual frames were obtained through the $`K^{}`$ filter and for field B, 11 such frames were obtained giving a total of 16 and 11 minutes, respectively, of on-target integration time. The mean HJD for the two datasets is 2449978.5198 and 2449978.5661 respectively.

## 3 Data reduction

### 3.1 Pre-processing

All the data were corrected for non-linearity using the ctio.irlincor command within IRAF111IRAF is distributed by the National Optical Astronomical Observatories, which is operated by AURA, Inc., under cooperative agreement with the NSF.. Polynomial coefficients of coeff1=1.0, coeff2$`=3.808477\times 10^2`$, coeff3$`=5.37237\times 10^2`$ were adopted on the basis of the reported non-linearity for the array (Lidman, Lidman95 (1995)).

Flatfields were produced as the difference between exposures of the white screen in the dome with the quartz lights on and off, respectively. The IRAC-2 camera is known to suffer from parasitic light in the flatfields at the 5% level, and a standard star was mapped in a 4 by 4 grid across the array on the night of the IC4499 observations as well as on the previous night in order to determine the necessary correction for the flat fields. The counts in each position were fitted by a two dimensional second order polynomium for each night. The results were in good agreement, to within 3%, and the averaged correction was applied to the flat field frame.

Sky frames were computed as the average of four dithered sky exposures obtained before and after the science frame. The “avsigclip” algorithm of the IRAF imcombine command with a threshold value of $`\sigma =5`$ was employed to eliminate stars in the sky fields. For each target frame the corresponding sky frame was subtracted and bad pixels were removed by linear interpolation. The resulting targets frames were finally averaged to produce the final target frames.

### 3.2 Standard star photometry

Fourteen standard star measurements were performed on the night Oct. 1/2, 1995, to determine the color terms of the camera in the $`J`$ and $`K^{}`$ filters and the airmass coefficients. Each star was observed 4 times with the star imaged in four different positions on the array. The stars were selected from the list of Carter and Meadows (Carter95 (1995)) to cover as large a range in $`(JK)`$ as possible, i.e. from $`0.02`$ to $`0.96`$, and they were observed at airmasses ranging from 1.03 to 1.80. Using the conversion from the Carter system to the CIT system determined by Carter (Carter93 (1993)) we find:

$`K_{\text{CIT}}`$ $`=`$ $`K_{\text{inst}}^{}+0.014(\pm 0.02)\times (JK^{})_{\text{inst}}`$
$`0.17(\pm 0.03)\times X+\text{const}`$
with an RMS scatter of 0.024 mag. We thus find no significant color term, and will assume a zero effect in the following.

On the night of the IC4499 observations we also observed the standard star HD136879 (Carter and Meadows Carter95 (1995)) before and after the observations of IC4499 in a pattern of four positions on the detector array. The airmass for the standard star measurements was 1.44 and 2.13 respectively, while the airmass for the IC4499 observations ranged from 1.85 to 2.05 (IC4499 is a very southern object). The standard star measurements were interpolated in airmass to provide the appropriate offset for the IC4499 data. Curves of growth were determined from artificial aperture photometry of the standard star and an aperture radius of 20 pixels (10 arcsec) was used to determine the offset to the standard system. From the RMS scatter for Eq.3.2 we estimate the error of the $`K`$-band zero-point to be 0.025 mag.

### 3.3 Photometry

The DAOPHOT package (Stetson Stetson87 (1987)) within IRAF was used to find the stars within the averaged frames and for performing photometry on them by point spread function (PSF) fitting. Due to the sparse nature of IC4499 crowding is not an issue for PSF fitting photometry.

The aperture corrections for the PSF photometry were determined from a number of well exposed and reasonably isolated stars in the IC4499 fields. The aperture photometry was performed after all the stars, apart from the stars selected for offset determination, had been subtracted by the IRAF/DAOPHOT substar routine. The resulting aperture corrections for five stars in each field agreed to better than 0.02 mag.

The RR Lyrae stars were identified from the finding chart of WN96. The final magnitudes on the CIT system are tabulated in Table 1 together with the periods and the type of pulsation as determined by WN96. The pulsation type is indicated by “ab” for fundamental mode pulsators, “c” for first overtone, “d” for double mode, and “e?” for possible second overtone pulsators. The estimated errors provided by DAOPHOT have been tabulated as well. Note that the errors on the magnitudes for stars in field B are significantly larger than in field A due to the shorter integration time.

The observed magnitudes have been collected at a specific time as already mentioned in Sect.2 so with an appropriate ephemeris for each star we can determine the phase of observation. Alistair Walker has kindly provided the full optical light curves from WN96 in machine readable format and using their period estimates we have determined the epoch of maximum $`V`$-light for each of the stars. On this basis we have determined the phase corresponding to the observed $`K`$-magnitude. We have then used the $`K`$-band template light curves from Jones, Carney and Fulbright (Jones96 (1996)) to compute the corresponding intensity averaged mean $`K`$-magnitude, $`K`$, based on the type of pulsation and the $`V`$-band amplitude as given by WN96. For the RRd stars we have not attempted a phase correction and for the putative RRe stars we have assumed that they are actually RRc stars as done by Castellani, Caputo and Castellani (Castellani03 (2003)). These estimates of the mean $`K`$-magnitude have also been tabulated in Table 1.

## 4 Results

### 4.1 The $`M_K`$-logP relation

The $`M_K`$\- $`\mathrm{log}P`$ relation has been established by Longmore et al. (Longmore90 (1990)) who determined the slope from photometry of RR Lyrae stars in 8 different galactic globular clusters. Recently, B03 has developed the relation further to incorporate a small but significant metallicity effect and they find the relation:

$$M_K=0.7702.101\mathrm{log}P+0.231\text{[Fe/H]}$$
(2)
For a globular cluster like IC4499 where we assume that there is no significant spread in \[Fe/H\] among the RR Lyrae stars, we can simply insert the best estimate of the metallicity in Eq.2 to get a $`M_K`$-$`\mathrm{log}P`$ relation. WN96 discusses the available metallicity estimates for IC4499 and find a value of $`\text{[Fe/H]}=1.65\pm 0.1`$ based on spectroscopic measurements of Ca II triplet lines in four RGB stars by Cannon (Cannon92 (1992)) and from a $`(V,(VI))`$ color-magnitude diagram.

We insert this value in Eq.2 and for each star we can then determine the distance modulus as:

$$(mM)_K=KM_K=K1.1512.101\mathrm{log}P_0$$
(3)
The relation refers to the fundamental period and thus it is necessary to fundamentalize the periods of the higher order pulsators. We adopt the period ratios as used by WN96, i.e. (P2:P1:P0)=(1.575:1.3439:1.0) and thus $`\mathrm{log}P_0=\mathrm{log}P_1+0.128`$ and $`\mathrm{log}P_0=\mathrm{log}P_2+0.197`$ for the RRc and RRe stars, respectively.

In Fig.1 we have plotted the $`K`$ magnitudes for the individual stars versus $`\mathrm{log}P_0`$, the fundamentalized period. Again we have assumed that the putative RRe stars are actually RRc stars. They actually fall exactly on the general relation. If we choose to fundamentalize the period of these stars as second overtone stars they would appear where the open diamonds in Fig.1 have been plotted which gives a less satisfying match.

Inspecting Fig.1 further, we see that three of the stars, marked with open symbols (V14, V32 and V81), are overly bright for their period when compared to the bulk of the stars. WN96 also found both V32 and V81 to be over-luminous in the optical, and they suspected both stars to have an unresolved red companion.

The detection limit is fainter than the faintest RR Lyrae star by about 1 mag in the shallow field (B) so the sample is complete and can not suffer from Malmquist bias.

### 4.2 The reddening and the distance modulus

For each star we can now use Eq.3 to determine the distance modulus and we can then compute the weighted mean of the resulting distances. If we use the sample of 29 RRab , RRc , and RRd stars excluding the three outliers we find $`(mM)_K=16.555\pm 0.016`$ with $`\sigma =0.087`$ mag. If we include V14 and V32 we find $`(mM)_K=16.526\pm 0.023`$ with $`\sigma =0.12`$ mag, which is only slightly smaller than the previous result. The three possible RRe stars do not affect the distance estimate significantly, independently of how we choose to fundamentalize them. Finally, if we choose to simply use the random phase $`K`$ magnitudes as a proxy for the mean $`K`$-magnitude instead of the template approach described above, we find $`(mM)_K=16.553\pm 0.019`$ with $`\sigma =0.104`$, basically an identical result but with a slightly larger uncertainty.

Due to the relative insensitivity of the $`M_K`$-$`\mathrm{log}P`$ relation to reddening, we can determine an independent value for the reddening towards IC4499 from the mean $`V`$-magnitude of the RR Lyrae stars as determined by WN96 and the $`M_V`$-\[Fe/H\] relation. For consistency we use the calibrations presented by B03 for both relations, and thus in addition to Eq. 2 we adopt

$$M_V=0.177(\pm 0.07)\text{[Fe/H]}+0.718(\pm 0.07)$$
(4)
which is valid for metal-poor stars ($`\text{[Fe/H]}<1.6`$).

WN96 find the mean $`V`$-magnitudes of the RR Lyrae stars to be $`<V_{\text{RR}}>=17.652\pm 0.006`$ and thus $`(mM)_V=17.226\pm 0.07`$. Adopting $`A_V=3.1\times E(BV)`$ and $`A_K=0.11\times A_V`$ (Cardelli et al. Cardelli89 (1989)) and using the constraint that the reddening corrected distance estimates from both the $`M_K`$-$`\mathrm{log}P`$ and $`M_V`$-\[Fe/H\] relations should be the same, we can write the relation:

$$(mM)_VA_V=(mM)_KA_K$$
(5)
from which follows:

$`A_V`$ $`=`$ $`[(mM)_V(mM)_K]/0.89`$ (6)
$`A_V`$ $`=`$ $`0.755\pm 0.08`$ (7)
$`E(BV)`$ $`=`$ $`0.24\pm 0.03`$ (8)
Our best estimate of the reddening is then $`E(BV)=0.24\pm 0.03`$ and the best estimate of the distance to IC4499 is $`(mM)_0=16.5550.11\times 0.755=16.472\pm 0.016`$ where the error estimate is the formal fitting error only. The DIRBE reddening map (Schlegel et al. Schlegel98 (1998)) gives $`E(BV)=0.23`$ at the position of IC4499 suggesting that the reddening is completely dominated by galactic reddening.

In addition to the fitting error we add the error contributions from the $`K`$-band zero-point (0.025 mag), from the aperture correction (0.02 mag), from the uncertainty in $`A_K`$ (0.009 mag), and from the uncertainty in the metallicity estimate ($`0.231\times 0.1=0.023`$ mag). This gives a final random error of 0.04 mag. The zero-point of the PL relation in Eq.2 has an uncertainty of 0.044 mag which should be considered as a systematic error. We should also keep in mind that the metallicity scale is not uniquely defined and that there might be systematic differences between scales reaching 0.2 dex (Layden Layden94 (1994)). This would constitute a systematic error of $`0.231\times 0.2=0.046`$ mag. This adds up to an overall systematic error estimate of 0.06 mag. A much larger potential systematic error arises from the adopted zero-point of the $`M_K`$-$`\mathrm{log}P`$ -\[Fe/H\] relation. As an example, Jones et al. (JCSL (1992)) found a zero-point based on Baade-Wesselink analysis of field RR Lyrae stars which is about 0.3 mag fainter than the B03 relation. A similar result was found by Layden et al. (Layden96 (1996)) from a statistical parallax study of field RR Lyrae stars. Walker (Walker92 (1992)) showed that adopting the Baade-Wesselink and statistical parallax zero-point would cause the LMC distance estimate from RR Lyrae stars to be about 0.3 mag shorter than the canonical Cepheid based distance of 18.5 mag. This shows that the B03 relation is in good agreement with the canonical Cepheid distance scale but it disagrees with the short distance scale. B03 also show that their relation is in excellent agreement with the HST parallax measurements of RR Lyr itself (Benedict et al. Benedict02 (2002)). Kovács Kovacs03 (2003) has redone the Baade-Wesselink analysis for a sample of galactic Cepheids similar to the one of Jones et al. (JCSL (1992)) using the same calibration as has been successfully used for Cepheids and he now finds agreement with the canonical LMC distance. We cannot explain the difference between the distance scales with the current data, but we simply choose to adopt the B03 zero-point to ensure that we are on the canonical Cepheid distance scale.

On this basis we find $`(mM)_0=16.47\pm 0.04`$ (random) $`\pm 0.06`$ (systematic), or $`d=19.7\pm 0.4\pm 0.6`$ kpc, as our best distance estimate to IC4499, which is in excellent agreement with the value of $`(mM)_0=16.45\pm 0.05`$ found by WN96. Our reddening value of $`E(BV)=0.24\pm 0.03`$ also agrees well with the value of $`E(BV)=0.22\pm 0.02`$ from WN96.

We stress that for stellar systems containing a reasonable number of RR Lyrae stars with known periods and metallicities, the $`M_K`$-$`\mathrm{log}P`$ -\[Fe/H\] relation combined with single epoch $`K`$-band data gives accurate and reddening insensitive distances for a very modest amount of observational effort as the dominating source of error is not the number of individual phase points.

### Acknowledgments

The data presented here was acquired during Directors Discretionary Time, granting of which is gratefully acknowledged. Thanks is due to Alistair Walker for providing a preprint of WN96 prior to publication, which helped the identification of the stars very much. Thanks is also due to Giuseppe Bono for comments on an early version of this paper and to an anonymous referee for suggestions which clarified the presentation.


text:
# Inflation and deflation in financial markets

## 1 Introduction

Financial markets have been receiving increasing attention from physicists \[1-10\]. One of the reasons for this is that recent empirical studies on financial time series have found power laws, which emerge from complex systems, for many of these series \[1-4\]. In the present study, we focus attention on inflationary bubbles which ended in crashes in financial markets. In the last two decades, the financial markets have been distinguished by increasingly frequent financial crashes. Unfortunately, although it is extremely important to elucidate the mechanisms of financial crashes, little is yet known . We here present the results of an empirical analysis performed by taking a new approach to bubbles and their collapses. We conducted a quantitative investigation on the statistical properties of price fluctuations in the Nikkei 225 index, which is the arithmetic average of the stock prices of 225 large companies listed in the Tokyo Stock Exchange. Figure 1(a) shows the daily series of the Nikkei 225 index during the period from January of 1975 to December of 2002. The Nikkei 225 index reached a high of almost 40,000 yen on the last trading day of the decade of the 1980s, but then declined approximately 14,309 during the period from the first trading day of 1990 to mid-August of 1992, a drop of about 63 percent. Thus, we divide the time series of the Nikkei 225 index into two parts: the period of inflation (bubbles), from January of 1975 to December of 1989, and the period of deflation (anti-bubbles), form January of 1990 to December of 2002. We show that the price fluctuations in the inflationary period can be approximated by asymptotic power laws provided by non-extensive statistical mechanics, while the price fluctuations in the deflationary period obey exponential laws provided by Boltzmann-Gibbs statistical mechanics.

## 2 Empirical analysis

### 2.1 Log returns

We first investigated the basic statistics of the price fluctuations. We define price fluctuations $`R_t`$ as the log return, $`R_t=\mathrm{ln}S_t\mathrm{ln}S_{t1}`$, where $`S_t`$ denotes the Nikkei 225 index on the date $`t`$. Figure 1(b) shows the daily returns of the Nikkei 225 index in the period from January 4, 1975 to December 30, 2002. Table 1 shows differences in the statistical properties between the inflationary period (January 4, 1975 to December 29, 1989) and the deflationary period (January 4, 1990 to December 30, 2002). The mean of the returns $`R`$ is positive in the inflationary period and negative in the deflationary period. The standard deviation of the returns $`\sigma `$ in the deflationary period is about two times as high as that in the inflationary period. Note that we omit the data of November 20, 1987, when the Nikkei 225 index had the greatest negative return, as an extreme value because extreme values would strongly affect the estimation of skewness $`s=E(R_tR)^3/\sigma ^3`$ and kurtosis $`k=E(R_tR)^4/\sigma ^4`$. The return distribution is negatively skewed during the inflationary period, and positively skewed in the deflationary period. Finally, the kurtosis of the return distribution in the inflationary period is higher than that in the deflationary period, indicating that the return distribution has a fat tail in the inflationary period, compared to that in the deflationary period. That is, there is a higher probability for extreme values in the return distribution in the inflationary period than in the deflationary period.

### 2.2 Absolute log-returns

For the purposes of the present study, we define volatility as the absolute value of the log return $`V_t=|R_t|`$. Figure 1(c) shows the absolute log-return $`V_t`$ for the Nikkei 225 index for the 28-year study period. Figures 2(a) and 2(b) show the semi-log plots of the complementary cumulative probability distributions of the absolute log-returns for the Nikkei 225 index for both the inflationary and deflationary periods. It follows from Figure 2(a) that the complementary cumulative probability distribution of the absolute log-returns for the period of deflation is accurately described by an exponential distribution,

$$P(V>x)\mathrm{exp}(\beta x),$$
(1)
in the whole range of the absolute log-returns $`V_t`$. We estimate the parameter $`\beta `$ using the least square method, which gives $`\beta =1.45`$ $`(0.0)`$. It also follows from Figure 2(b) that the complementary cumulative probability distribution of the absolute log-returns for the period of inflation deviates from the exponential distribution as the absolute log-returns becomes greater. Figure 3(a) shows the log-log plot of the distribution of the absolute log-returns in the inflationary period. This plot shows an asymptotic power law where the value of the absolute log-return is large, but deviates from the power law where the value of the absolute log-return is small. One approach to incorporating this deviation from the power law is to consider the $`q`$-exponential function . We find that the complementary cumulative probability distribution of the absolute log-returns of the inflationary period is well approximated by the following $`q`$-exponential distribution:

$$P(V>x)[1(1q)\beta x]^{1/(1q)}.$$
(2)
We estimate the parameters, $`\beta `$ and $`q`$ using the nonlinear least square method, which gives $`\beta =2`$ and $`q=1.14`$. Note that for $`q>1`$ the $`q`$-exponential distribution has a power-law tail, i.e., $`P(V>x)x^\alpha `$ , and for $`q=1`$ it becomes an exponential function since $`lim_{q1}\mathrm{exp}_q(\beta x)=\mathrm{exp}(\beta x)`$. Here, we introduce another possible way of analyzing data by using a generalized mono-log plot based on the inverse function of the $`q`$-exponential function, which is the $`q`$-logarithmic function given by $`\mathrm{ln}_q[P(V>x)](x^{1q}1)/(1q)`$. This generalized function arises naturally in the framework of Tsallis statistics . It is easy to verify that the plot of $`\mathrm{ln}q[r(x)]`$ versus $`x`$ leads to a straight line if the data are well described by the $`q`$-exponential distribution (Figure 3(b)).

## 3 Concluding remarks

Our results indicate that the absolute log-return is approximated by the $`q`$-exponential distribution where $`q=1.14`$ in the inflationary period, while during the deflationary period, the distribution of the absolute log-returns is accurately described by an exponential distribution, that is, the $`q`$-exponential distribution where $`q1`$. As has been demonstrated by C. Tsallis , Boltzmann-Gibbs statistical mechanics typically provides exponential laws for describing stationary states, while non-extensive statistical mechanics typically provides asymptotic power laws. The empirical findings presented here indicate that the price fluctuations in the inflationary (bubble) period can be described by a power law based on non-extensive statistical mechanics, while the price fluctuations in the deflationary (anti-bubble) period can be described by an exponential law based on Boltzmann-Gibbs statistical mechanics. The present study raises the interesting question, left for future research, of whether the empirical facts found here hold true in other stock markets such as the New York Stock Exchange. This study will be left for future work. Our empirical results suggest that the momentous structural change was cased and the $`q`$-value changed from greater than unity to unity at the beginning of 1990 when the speculative bubble collapsed in Japan’s stock markets. However, to date, no model so far has successfully explained such changes in the statistical properties of price fluctuations. Hence the next step is to model such behavior in stock markets.

## 4 Acknowledgements

The author would like to thank Prof. Thomas Lux for valuable comments and criticism, and Ms. Michiyo Kaizoji for her assistance in collecting the data. Financial support from the Japan Society for the Promotion of Science under Grant-in-Aid No. 06632 is gratefully acknowledged. All remaining errors, of course, are mine.