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# Invariant Properties of the Ansatz of the Hirota Method for Quasilinear Parabolic equations

## 1

We consider the quasilinear parabolic equation

$`(1+b_1u+b_2u^2)u_t(h_1+h_2u+h_3u^2)u_xg_0(u_x)^2`$
$`(k_0+k_1u)u_{xx}+{\displaystyle \underset{i=1}{\overset{4}{}}}\phi _iu^i=0.`$ (1.1)
where $`b_j(x,t)`$, $`h_j(x,t)`$, $`k_{j1}(x,t)`$, $`j=1,2`$, $`\phi _i(x,t)`$, $`i=1,\mathrm{},4`$, $`g_0(x,t)`$ are smooth functions or constants.

It is assumed that the coefficients of $`u_t`$ and$`u_{xx}`$ do not vanish in the range of the function $`u`$.

We seek a solution in the form of the fraction  (R. Hirota):

$$u(x,t)=\frac{G(x,t)}{F(x,t)}.$$
(1.2)
After multiplication by $`F(x,t)^4`$, for the functions $`G(x,t)`$ and $`F(x,t)`$ we receive a homogeneous equation of the fourth order of homogeneity with respect to all functions.

A special case of Eq. (1.1) has the form

$$(1+b_1u)u_t(h_1+h_2u)u_xk_0u_{xx}+\underset{i=1}{\overset{3}{}}\phi _iu^i=0,$$
(1.3)
which, after the substitution of (1.2) and multiplication by $`F(x,t)^3`$ becomes an equation of the third order of homogeneity with respect to the functions $`G(x,t)`$ and $`F(x,t)`$.

Thus this equation is of the third order of homogeneity, and the the components $`u^4`$ or $`(u_x)^2`$ complicate the equation so that it is an equation of the fourth order of homogeneity.

Special cases of this equation are widely known (see the extensive bibliography in ).

If $`h_1=0`$, $`\phi _i=0`$, $`i=1,3`$, then this is the Bürgers equation in the wave theory .

If $`h_1=0`$, $`h_2=0`$, $`\phi _1=1`$, $`\phi _2=\phi _1`$, $`\phi _3=0`$, then this is the Fisher–Kolmogorov–Petrovskii–Piskynov (FKPP) equation (see the bibliography in );

If $`h_1=0`$, $`h_2=0`$, $`\phi _1=1`$, $`\phi _2=0`$, $`\phi _3=\phi _1`$, then this is the FitzHygh–Nagumo–Semenov (FhNS) equation in the wave theory (see bibliography in ), , p. 12), and the Allen–Cahn equation (which is close to the FhNS equation) in the theory of phase transitions  and in ecology . Equations in which the coefficients of the Laplace operator depend on the function $`u`$ appear in models of the theory of epidemic distribution . Equations in which the coefficients of the derivative $`u_t`$ depending on $`u_x`$ arise, for example, in models of processes of ocean water freshening . The initial reasons on a theme of the given operation can be detected in operation .

The method proposed for constructing solutions of Eq. (1.1) ideologically goes back to the Hirota method. An essential distinguishing feature of the present paper is that a new invariance property of solutions of equations from a given class is discovered (cf. ). The equations obtained from (1.1) after the substitution of (1.2) can be arranged in the order of homogeneity as embedded equations.

From now on, we use the following notation for the derivatives:

$$_\nu ()=()_\nu ,\nu =x,t$$
(1.4)
We choose the solution (1.2) of Eq. (1.3) in the form

$`G(x,t)`$ $`=`$ $`\mathrm{exp}(H(x,t))B(x,t)+Q(x,t),`$ (1.5)
$`F(x,t)`$ $`=`$ $`\left(1+A(x,t)\mathrm{exp}(H(x,t))\right).`$ (1.6)
All the other representations of the solution obtained by the Hirota method can be reduced to this form. There is only one possibility to extend this notation, namely, to add sums of series in powers of the exponentials $`\mathrm{exp}(iH(x,t))`$, $`i=2,\mathrm{}`$. Let us substitute (1.5) first into (1.2) and then into (1.3) and equate the coefficients of equal powers of $`\mathrm{exp}(H(x,t))`$ (these powers are $`0,1,2,3`$) with zero.

We denote the obtained system of equations by $`(eq0)`$, $`(eq1)`$, $`(eq2)`$, $`(eq3)`$.

Equation $`(eq0)`$ is the same original equation (1.3) for the function $`Q(x,t)`$. Hence $`Q(x,t)`$ is a solution of this equation.

Let us consider the equation obtained at a power of $`\mathrm{exp}(3H(x,t))`$, namely, Eq. $`(eq3)`$:

$`B\left(B^2\phi _32k_0(A_x)^2+B(b_1A_t+H_2A_x)\right)`$
$`+A\left(B^2(\phi _1+\phi _3)+2k_0A_xB_x+B(A_t+B_1B_t+h_1A_xh_2B_x+k_0A_{xx})\right)`$
$`+A^2(B\phi _1+B_th_1B_xK_0B_{xx})=0.`$ (1.7)
The special case considered in , p. 51, Section 24, and in , p. 190, is the case of $`B=(\pm )A`$, where the equation has the form

$$A^2B+B^3ABA_t2BA_x^2+ABA_{xx}+A^2B_t+2AA_xB_xA^2B_{xx}=0.$$
(1.8)
and is an identity.

First, we present a lemma in which we prove that, in the conventional ansatz form of a solution, there is considerable arbitrariness related to the “extra” function $`A(x,t)`$, which must be excluded. Namely, this function must be added to the function $`H(x,t)`$ according to the formula $`H(x,t)_{(\text{new})}=H(x,t)_{(\text{old})}+\mathrm{exp}(\mathrm{ln}(A))`$. Then the same notation $`H(x,t))_{(\text{new})}`$ is preserved for the function $`H(x,t)`$.

###### Lemma 1.1

Suppose that the functions $`B(x,t)`$, $`A(x,t)`$ in (1.5) belong to $`C^2[R^1R_+^1]`$ and satisfy Eq. $`(eq3)`$, and the function $`B(x,t)`$ has the form

$$B(x,t)=M(x,t)A(x,t).$$
(1.9)
Then the function $`M(x,t)`$ satisfies Eq. (1.3)

$$\left(1+b_1(x,t)M\right)M_t\left(h_1(x,t)+h_2(x,t)M\right)M_xk_0(x,t)M_{xx}+\underset{i=1}{\overset{3}{}}\phi _i(x,t)M^i=0$$
(1.10)
for all functions $`A(x,t)`$.

Proof. The proof of the lemma is directly obtained by the substitution of (1.9) into (1.7). Equation (1.10), which follows from Eq. (1.7), coincides with (1.3). Hence the function $`M(x,t)`$ is a solution of this equation. Below $`M(x,t)`$ stands for some known exact solution of Eq. (1.3). $`\mathrm{}`$

It follows from Lemma 1.1 that the number of arbitrary functions in the ansatz can be decreased. Namely, we write $`A=\mathrm{exp}(\mathrm{ln}(A))`$ and add this summand to the function $`H(x,t)`$. In the present paper, solutions of the form (1.2) are called solutions of the first, second, etc. mode with respect to their complexity (the number of terms containing a power of $`\mathrm{exp}(iH)`$ as a factor).

###### Definition 1.1

The set $`\mathrm{\Omega }`$ of solutions of Eq. (1.3) of the form

$$u(x,t)=\frac{\mathrm{exp}(H(x,t))M(x,t)+Q(x,t)}{1+\mathrm{exp}(H(x,t))},$$
(1.11)
where $`M(x,t)`$, $`Q(x,t)`$ are some solutions of Eqs. (1.3) and $`H(x,t)`$ is an arbitrary function, will be called SOLUTIONS OF THE FIRST MODE; solutions of the form

$$U(x,t)=\frac{\mathrm{exp}(H(x,t))M(x,t)+\mathrm{exp}(2H(x,t))R(x,t)+Q(x,t)}{1+Z_1(x,t)\mathrm{exp}(H(x,t))+\mathrm{exp}(2H(x,t))Z_2(x,t)},$$
(1.12)
will be called SOLUTIONS OF THE SECOND MODE; etc.

We have the following theorem.

###### Theorem 1.1

Suppose that all coefficients in Eq. (1.3) are constants. The function $`M`$ is a solution of the equation

$$(1+b_1M)M_t(h_1+h_2M)M_xk_0M_{xx}+\underset{i=1}{\overset{3}{}}\phi _iM^i=0,$$
(1.13)
and $`Q`$ is a solution of the equation

$$(1+b_1Q)Q_t(h_1+h_2Q)Q_xk_0Q_{xx}+\underset{i=1}{\overset{3}{}}\phi _iQ^i=0,$$
(1.14)
In addition, they are related by the compatibility relations

$`((MQ)b_1k_0H_{xx}(x,t)`$
$`+H_x(x,t)\left(M(b_1h_1h_2)+Q(b_1h_1+h_2)+2b_1k_0(M_xQ_x)\right)`$
$`+(2+Mb_1+Qb_1)k_0H_x(x,t)^2`$
$`+(MQ)(M^2b_1\phi _3Q^2b_1\phi _3+M(\phi _3+b_1(\phi _1+\phi _3))Q(\phi _3+b_1(\phi _1+\phi _3))`$
$`+b_1(b_1(M_tQ_t)+h_2(M_x+Q_x))))=0,`$ (1.15)
and, together with the function $`H(x,t)`$ satisfy the equation

$$(MQ)b_1H_t(x,t)+H_x(x,t)(Mh_2Qh_22k_0H_x(x,t))+(MQ)^2\phi _3=0.$$
(1.16)
Then the function $`u(x,t)`$ is a solution of Eq. (1.3).

(The arguments of the functions are omitted here for brevity.).

###### Remark 1.1

The equation (1.13) is the Riccaty equation for the functions $`H^{(1,0)}(x,t)`$, and equation (1.14) is the Hamilton–Jacobi equation for the function $`H(x,t)`$.

If $`b_1=0`$, then the statement of Theorem 1.1 becomes simpler and the function $`H(x,t)`$ can explicitly be written in terms of the functions $`M(x,t)`$ and $`Q(x,t)`$.

###### Theorem 1.2

Let $`b_1=0`$. Suppose that all coefficients in Eq. (1.3) are constants. The function $`M`$ is a solution of the equation

$$M_t(h_1+h_2M)M_xk_0M_{xx}+\underset{i=1}{\overset{3}{}}\phi _iM^i=0,$$
(1.17)
and $`Q`$ is a solution of the equation

$$Q_t(h_1+h_2Q)Q_xk_0Q_{xx}+\underset{i=1}{\overset{3}{}}\phi _iQ^i=0.$$
(1.18)
In addition, they are related by the compatibility relations

$`\left({\displaystyle (M_tQ_t)𝑑x}\right)4(h_2m)2(M^2Q^2)(h_2^2+12k_0\phi _3h_2m)`$
$`4(MQ)(h_1h_2+4k_0\phi _2h_1m)+16k_0C_0^{}(t)`$
$`+4k_0(h_2+3m)(M_xQ_x)=0,`$ (1.19)
where $`m=((\pm )\sqrt{(h_2^2+8k_0\phi _3)})`$, and the function $`H(x,t)`$ has the form

$$H(x,t)=C_0(t)\frac{(M(x,t)Q(x,t))\left(h_2\pm \sqrt{h_2^2+8k_0\phi _3}\right)}{4k_0}𝑑x.$$
(1.20)
Then the function $`u(x,t)`$ is a solution of Eq. (1.3) .

Proof. We give a brief proof of this theorem. We substitute (1.11) into (1.3), reduce to the common denominator, and consider the numerator of the obtained large equation.

We equate with zero the coefficient of the exponential raised to the zero power, i.e., the coefficient of $`\mathrm{exp}(0)`$. The equation thus obtained is exactly Eq. (1.3) for the function $`Q(x,t)`$. Hence the function $`Q(x,t)`$ is a solution of this equation.

We equate with zero the coefficient of the exponential $`\mathrm{exp}(3H(x,t))`$. The equation thus obtained will be called Eq. $`(eq3)`$. Again it is exactly Eq. (1.3). Hence the function $`M(x,t)`$ is a solution of this equation.

We equate with zero the coefficient of the exponential $`\mathrm{exp}(H(x,t))`$. The equation thus obtained will be called Eq. $`(eq1)`$.

We equate with zero the coefficient of the exponential $`\mathrm{exp}(2H(x,t))`$. The equation thus obtained will be called Eq. $`(eq2)`$. Since the functions $`Q`$ and $`M`$ satisfy Eq. (1.3), we exclude the second derivatives of the functions $`Q_{xx}`$ and $`M_{xx}`$ from Eqs. $`(eq1)`$ and $`(eq2)`$.

We express $`H_{xx}`$ from Eq. $`(eq1)`$ and stubstitute it into Eq. $`(eq2)`$. Then we obtain the Hamilton–Jacobi equation (1.16). For $`b_1=0`$ this equation in Theorem 1.2 naturally determines the function $`H(x,t)`$ (1.20). If we express the derivative $`H_t`$ from Eq. (1.16) and substitute it into Eq. $`(eq1)`$, then we obtain the compatibility condition (1.15) or (1.19), respectively. If we substitute the derivative $`H_t`$ into Eq. $`(eq2)`$, then we also obtain the compatibility condition (i.e., the two equations coincide). In conclusion we note that there exist two branches in the relations given in the theorem with different signs. The proof of the theorem is complete. $`\mathrm{}`$

In fact, the theorems proved above describe the procedure of constructing new solutions of the equations (an analog of the superposition of solutions), which is specified by the properties of Eq. (1.3) and the ansatz.

Let us consider two smooth solutions of Eq. (1.3): $`M(x,t)`$ and $`Q(x,t)`$. The following function is assign to these solutions:

$$P(Q,M)=\frac{Q(x,t)+M(x,t)\mathrm{exp}(H(x,t))}{1+\mathrm{exp}(H(x,t))},$$
(1.21)
where $`H(x,t)`$ is determined by (1.16) or by formula (1.20) if $`b_1=0`$.

The following theorem states that the solutions are invariant.

###### Theorem 1.3

Suppose that two pairs of functions

$$(M_0(x,t),Q_0(x,t))\text{and}(M_1(x,t),Q_1(x,t))$$

are solutions of the equations

$$(1+b_1M_j)M_{j}^{}{}_{t}{}^{}(h_1+h_2M_j)M_{j}^{}{}_{x}{}^{}k_0M_{j}^{}{}_{xx}{}^{}+\underset{i=1}{\overset{3}{}}\phi _iM_{j}^{}{}_{}{}^{i}=0,$$
(1.22)
$$(1+b_1Q_j)Q_{j}^{}{}_{t}{}^{}(h_1+h_2Q_j)Q_{j}^{}{}_{x}{}^{}k_0Q_{j}^{}{}_{xx}{}^{}+\underset{i=1}{\overset{3}{}}\phi _iQ_{j}^{}{}_{}{}^{i}=0,$$
(1.23)
where $`j=0,1`$, and that they determine the solution by formula (1.11). Suppose that the assumptions of Theorems 1.1 and 1.2 hold for each pair. Then the two functions

$`Q(x,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(H_1(x,t))M_1(x,t)+Q_1(x,t)}{1+\mathrm{exp}(H_1(x,t))}},`$
$`M(x,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(H_0(x,t))M_0(x,t)+Q_0(x,t)}{1+\mathrm{exp}(H_0(x,t))}}`$ (1.24)
are solutions of the equations

$$(1+b_1M)M_t(h_1+h_2M)M_xk_0M_{xx}+\underset{i=1}{\overset{3}{}}\phi _iM^i=0,$$
(1.25)
$$(1+b_1Q)Q_t(h_1+h_2Q)Q_xk_0Q_{xx}+\underset{i=1}{\overset{3}{}}\phi _iQ^i=0,$$
(1.26)
and the function

$$u(x,t)=\frac{\mathrm{exp}(H(x,t))M(x,t)+Q(x,t)}{1+\mathrm{exp}(H(x,t))}$$
(1.27)
is a solution of Eq. (1.3).

Proof. Each of the functions $`M_j`$, $`Q_j`$, $`j=0,1,`$, satisfies Eq. (1.3) written for this function. Pariwise, they satisfy the corresponding compatibility condition (1.15) or (1.2). The theorem can be proved by a direct substitution. $`\mathrm{}`$

###### Remark 1.2

After calculating the derivative with respect to $`x`$, it is possible to consider the compatibility condition (1.15) as a modified Bürgers equation with a potential.

###### Corollary 1.1

The class of solutions (1.27) is closed with respect to the procedure described in Theorem 1.3. Thus (1.27) defines a surjective mapping:

$$\text{Vol}:\mathrm{\Omega }\mathrm{\Omega }.$$
(1.28)
Similar results take place for Eq. (1.1). We perform the change of the variables $`\phi _1=a_2a_1`$, $`\phi _2=(a_1+a_2+a_1a_2)`$, $`\phi _3=1+a_1+a_2`$, $`\phi _4=1`$, which is equivalent to a source-sink function written so that the four roots are explicitly distinguished: $`u(a_1u)(a_2u)(1u)`$.

We have a lemma similar to Lemma 1.1 and the following theorem.

###### Theorem 1.4

Suppose that the solution (1.1) has the form (1.11) and the relations hold in one of the following cases:

1) $`(Q(x,t)`$ or $`M(x,t))const=k_0/k_1`$, $`(k_0+k_1)(k_0+a_1k_1)(k_0+2a_1k_1)=0`$;

2) $`b_1=k_1/k_0`$, $`b_2=0`$;

3) $`b_1=0`$, $`b_2=0`$, $`k_1=0`$.

Then the functions $`M`$ and $`Q`$ are also solutions of an equation coinciding with Eq. (1.1). Moreover, these functions are related by the compatibility relations (the Riccati equation for $`H_x(x,t)`$). In addition, the Hamilton–Jacobi equation must be satisfied.

(Here we do not write these equations, since they are very cumbersome.) The proof is similar to the proofs of the preceding theorems.

## 2

In the case of variable coefficients, there are additional possibilities to satisfy the system of equations. Since we construct exact solutions, not all coefficients of the equation are arbitrary.

We have the following theorem.

###### Theorem 2.1

Suppose that $`b_1(x,t)=0`$, the function $`M(x,t)`$ is a solution of the equation

$$M_t\left(h_1(x,t)+h_2(x,t)M\right)M_xk_0(x,t)M_{xx}+\underset{i=1}{\overset{3}{}}\phi _i(x,t)M^i=0,$$
(2.29)
the function $`Q(x,t)`$ satisfies the modified equation

$$Q_t\left(h_{01}(x,t)+h_{02}(x,t)Q\right)Q_xk_{00}(x,t)Q_{xx}+\underset{i=4}{\overset{6}{}}\phi _i(x,t)Q^{i3}=0,$$
(2.30)
and the coefficients in the equations satisfy the relations

$`Q\left(\phi _1(x,t)\phi _4(x,t)+Q(\phi _2(x,t)\phi _5(x,t))+Q^2(\phi _3(x,t)\phi _6(x,t))\right)`$
$`+\left(h_{01}(x,t)h_1(x,t)+(h_{02}(x,t)h_2(x,t))Q\right)Q_x`$
$`(k_0Q_{xx}k_{00}Q_{xx})=0,`$ (2.31)
and the functions $`M`$, $`Q`$ and the coefficients satisfy the second compatibility equation

$`\left({\displaystyle _t\left(\frac{(M+Q)(h_2\pm m)}{k_0}\right)dx}\right)4k_0(M+Q)m`$
$`+2(MQ)(\pm h_2^3(M^2Q^2)+2h_1(MQ)((\pm )m^2+h_2m)`$
$`+h_{2}^{}{}_{}{}^{2}\left((MQ)(M+Q)m(\pm )2Mk_{0}^{}{}_{x}{}^{}+(\pm )2Qk_{0}^{}{}_{x}{}^{}+(\pm )6k_0(M_xQ_x)\right)`$
$`+2h_2((\pm )4k_0\phi _3(M^2Q^2)+(\pm )k_0h_{2}^{}{}_{x}{}^{}(MQ)m(k_{0}^{}{}_{x}{}^{}(MQ)+k_0(M_xQ_x))`$
$`+2k_0(m(2(MQ)(2\phi _2+3\phi _3(M+Q))4C_0^{}(t)+h_{2}^{}{}_{x}{}^{}(MQ)`$
$`(\pm )4\phi _3k_{0}^{}{}_{x}{}^{}(MQ)+(\pm )24\phi _3k_0(M_xQ_x)+(\pm )4\phi _{3}^{}{}_{x}{}^{}k_0(MQ)))=0,`$ (2.32)
where $`m=(\sqrt{(h_2^2+8k_0\phi _3)})`$ and the function $`H(x,t)`$ has the form

$$H(x,t)=C_0(t)\frac{(M(x,t)Q(x,t))\left(h_2(x,t)\pm \sqrt{h_2(x,t)^2+8k_0(x,t)\phi _3(x,t)}\right)}{4k_0(x,t)}𝑑x.$$
(2.33)
Then the function $`u(x,t)`$ (1.11) is a solution of Eq. (1.3).

Proof. The proof is similar to that of Theorem 1.2. Note that, first, in this case the set of solutions to the system of equations for the functions $`M`$, $`Q`$ and the coefficients of the equation described in the theorem is not empty;

second, it is possible to construct solutions with variable roots of the source-sink function.

We substitute (1.11) into (1.3), reduce to the common denominator, and consider the numerator of the obtained equation. By equating with zero the coefficient of the exponential raised to the zero power, we obtain Eq. $`(eq0)`$, which is the compatibility condition for coefficients (2.22). From this condition we express some coefficient and exclude this coefficient from the other equations. The equation obtained by equating the coefficient of the exponential $`\mathrm{exp}(H(x,t))`$ with zero is called $`(eq1)`$. The equation obtained by equating the coefficient of the exponential $`\mathrm{exp}(2H(x,t))`$ with zero is called $`(eq2)`$. Since the function $`Q`$ satisfies Eq. (2.30), we exclude the derivative $`Q_t(x,t)`$ from Eqs. $`(eq1)`$ and $`(eq2)`$. We can express $`H_{xx}(x,t)`$ from Eq. $`(eq1)`$ and substitute it into Eq. $`(eq2)`$. Then we obtain the Hamilton–Jacobi equation, which naturally determines the function $`H(x,t)`$ (2.24). The proof of the theorem is complete.

In the construction of the solution of the second mode, the transformation relates the solution of the original equation and the solutions of two modified equations, and so on.

The FitzHygh–Nagumo–Semenov equation

EXAMPLE 1. We apply our method to the FitzHygh–Nagumo–Semenov (FhNS) equation known in wave theory:

$$u_t\phi _3u_{xx}/(2a^2)\phi _3u+\phi _3u^3.$$
(2.34)
where $`k_0=\phi _3/(2a^2)`$ and $`\phi _1=\phi _3`$ are constants. We show that this method leads to new results and provides a new understanding of the results (see , p. 64, and ).

First we show the action of the semigroup related to the transformation (1.11) and explain how the translation constants pass into new solutions, by choosing the form of the functions (of its solutions, , p. 17, , p. 51) satisfying the assumptions of Theorem 1.2. Namely, we consider the functions $`u(x,t)=M(x,t)`$ or $`u(x,t)=Q(x,t)`$:

$`M(x,t)`$ $`=`$ $`(1+\mathrm{exp}(m_0+ax+bt))^1,`$ (2.35)
$`Q(x,t)`$ $`=`$ $`(1+\mathrm{exp}(2ax+q_0))/(1+\mathrm{exp}(2ax+q_0)),b=3\phi _3/2.`$
We calculate the function

$$H(x,t)=\mathrm{ln}\frac{1+\mathrm{exp}(m_0+ax3t\phi _3/2)}{1+\mathrm{exp}(q_0+2ax)},$$
(2.36)
and obtain the solution

$$u_1=\frac{1\mathrm{exp}(q_02ax+\mathrm{ln}(2))}{1+\mathrm{exp}(q_02ax+\mathrm{ln}(2))+\mathrm{exp}(m_0q_0ax3t\phi _3/2)}.$$
(2.37)
Thus the original transition constants $`q_0`$, $`m_0`$ pass into a new solution and a new constant $`\mathrm{ln}(2)`$ appears.

Note that this is a well-known solution that describes the wave interaction (, p. 190, , p. 51, and ). Below we prove that this solution can be expressed via the solution of the linear parabolic equation.

There exists another version of the functions for which the assumptions of Theorem 1.2 are satisfied:

$`M(x,t)`$ $`=`$ $`(1+\mathrm{exp}(m_1ax+bt))^1,`$ (2.38)
$`Q(x,t)`$ $`=`$ $`(1\mathrm{exp}(q_12ax))/(1+\mathrm{exp}(q_12ax)),b=3\phi _3/2.`$
Then the function $`H(x,t)`$ has the form

$$H(x,t)=\mathrm{ln}\frac{1+\mathrm{exp}(m_1ax3t\phi _3/2}{1+\mathrm{exp}(q_12ax)},$$
(2.39)
and we obtain the solution

$$u_2=\frac{1\mathrm{exp}(q_1+2ax+\mathrm{ln}(2))}{1+\mathrm{exp}(q_1+2ax+\mathrm{ln}(2))+\mathrm{exp}(m_1q_1+ax3t\phi _3/2)}.$$
(2.40)
For the new functions $`Q(x,t)=u_1`$, $`M(x,t)=u_2`$, we take the solutions of Eq. (2.34):

$`H(x,t)`$ $`=`$ $`\mathrm{ln}({\displaystyle \frac{2\mathrm{exp}(q_1)+\mathrm{exp}(2ax)+\mathrm{exp}(m_1q_1ax3t\phi _3/2)}{1+2\mathrm{exp}(q_02ax)}}`$
$`+\mathrm{exp}(m_0q_0ax3t\phi _3/2)),`$
$`u_3`$ $`=`$ $`\left(1+\mathrm{exp}\left(q_0q_1+2ax+\mathrm{ln}{\displaystyle \frac{2+\mathrm{exp}(q_1)}{2+\mathrm{exp}(q_0)}}\right)\right)`$ (2.41)
$`\times (1+\mathrm{exp}(q_0q_1+2ax+\mathrm{ln}{\displaystyle \frac{2+\mathrm{exp}(q_1)}{2+\mathrm{exp}(q_0)}})`$
$`+\mathrm{exp}(q_1+ax+\mathrm{ln}{\displaystyle \frac{\mathrm{exp}(m_1+q_0)+\mathrm{exp}(m_0+q_1)}{2+\mathrm{exp}(q_0)}}3{\displaystyle \frac{t\phi _3}{2}}))^1.`$
By analyzing this formula, it is possible to understand the structure of the translation constant. Indeed, in the variables $`\tau _1=ax+3t\phi _3/2`$, $`\tau _2=ax+3t\phi _3/2`$, the solution (2.41) has the form

$`u_3(\tau _1,\tau _2)`$ $`=`$ $`{\displaystyle \frac{c_1\mathrm{exp}(\tau _1)+c_2\mathrm{exp}(\tau _2)}{1+c_1\mathrm{exp}(\tau _1)+c_2\mathrm{exp}(\tau _2)}},`$
$`c_1`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(q_1)(2+\mathrm{exp}(q_0))}{\mathrm{exp}(q_0+m_1)+\mathrm{exp}(q_1+m_0)}},`$ (2.42)
$`c_2`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(q_0)(2+\mathrm{exp}(q_1))}{\mathrm{exp}(q_0+m_1)+\mathrm{exp}(q_1+m_0)}}.`$
###### Lemma 2.1

The set $`u(x,t)`$, $`M(x,t)`$, $`Q(x,t)`$ of solutions of the FhNS equation (2.34) of the form (2.38), (2.42) forms a semigroup with operation $``$ (1.11) and the following properties:

(a) commutativity, $`M(x,t)Q(x,t)=Q(x,t)M(x,t)`$;

(b) associativity, $`u(x,t)(M(x,t)Q(x,t))=(u(x,t)Q(x,t))M(x,t)`$;

(c) any element is the unity for itself, $`M(x,t)M(x,t)=M(x,t)`$.

Proof. Let us consider two solutions with different sets of constants $`C_i`$, $`V_i`$, $`i=1,2`$

$`Q(x,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(\tau _1+C_1)+\mathrm{exp}(\tau _2+C_2)}{1+\mathrm{exp}(\tau _1+C_1)+\mathrm{exp}(\tau _2+C_2)}},\tau _i=\tau _i[x,t],`$
$`M(x,t)`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(\tau _1+V_1)+\mathrm{exp}(\tau _2+V_2)}{1+\mathrm{exp}(\tau _1+V_1)+\mathrm{exp}(\tau _2+V_2)}},\tau _i=\tau _i(x,t).`$ (2.43)
By using Theorem 1.2, we find the solution

$`u(x,t)`$ $`=`$ $`M(x,t)Q(x,t)`$
$`=`$ $`{\displaystyle \frac{\mathrm{exp}(\tau _1)(\mathrm{exp}(C_1)+\mathrm{exp}(V_1))+\mathrm{exp}(\tau _2)(\mathrm{exp}(C_2)+\mathrm{exp}(V_2))}{1+\mathrm{exp}(\tau _1)(\mathrm{exp}(C_1)+\mathrm{exp}(V_1))+\mathrm{exp}(\tau _2)(\mathrm{exp}(C_2)+\mathrm{exp}(V_2))}},`$
$`\tau _i`$ $`=`$ $`\tau _i[x,t],`$ (2.44)
which implies the properties given in Lemma 2.1.

Note that the action of the semigroup is a nonlinear “time-translation” of the solution.

As was noted in , p. 192, , p. 64, , and , it is known that Eq. (2.34) has solutions with singularities (monsters, contrast structures) whose role cannot be explained.

We have the following theorem.

###### Theorem 2.2

Suppose that the function $`Q(x,t)`$ is a solution of the equation

$$Q_tk_0Q_{xx}\phi _1Q+\phi _3Q^3=0,$$
(2.45)
the function $`U(x,t)`$ is a solution of the linear parabolic equation

$$U_t+\phi _3a_1/a^2U_x\phi _3/a^2U_{xx}=0,$$
(2.46)
and the compatibility condition

$$Q_x+Q(a_1(\pm )aQU_x/U)=0$$
(2.47)
is satisfied. Then the function $`u(x,t)`$ is a solution of Eq. (2.34) and has the form

$$u=\frac{Q(x,t)U(x,t)\mathrm{exp}(t\phi _3(1+a_1^2/(2a^2)))}{Q(x,t)\mathrm{exp}(a_1x)U(x,t)\mathrm{exp}(t\phi _3(1+a_1^2/(2a^2)))}.$$
(2.48)
Proof. To prove this theorem, we substitute the ansatz of the solution of the third mode

$$u(x,t)=\frac{\mathrm{exp}(H(x,t))M(x,t)+\mathrm{exp}(2H(x,t))R(x,t)Q(x,t)}{1+Z(x,t)\mathrm{exp}(H(x,t))+\mathrm{exp}(2H(x,t))Z_1(x,t)+\mathrm{exp}(3H(x,t))R(x,t)},$$
(2.49)
into Eq. (2.34), collect similar terms with equal powers of the exponentials, equate them with zero, and obtain a system of nine equations. As previously, from these equations we successively exclude the second derivatives of all functions contained there. Since the ansatz is invariant, there is a remarkable fact common for the algorithm of constructing the solution. Namely, just as in the theorems presented above, the equation becomes factored at some step. This equaton contains various hints, which allow us to correct the ansatz (2.50) and to pass to the next iteration.

Equation $`(eq8)`$ has the form

$$2a^2\phi _3+2a^2H_t+\phi _3H_{x}^{}{}_{}{}^{2}\phi _3H_{xx}=0.$$
(2.50)
Its solution has the the form

$$H(x,t)=\mathrm{ln}(U(x,t))+a_1x(2a^2+a_1^2)t\phi _3/(2a^2),$$
(2.51)
and thus implies (2.47).

Equation $`(eq0)`$ is exactly Eq. (2.34) for the function $`Q(x,t)`$. Hence the function $`Q(x,t)`$ is a solution of this equation.

Equation $`(eq1)`$ at the second iteration step implies the compatibility conditions.

The second iteration of Eq. (eq1) implies the compatibility condition in the form of the Riccaty equation, in which we study only the situation corresponding to the sign “plus”.

Further, at each successive iteration, we have some versions dur to the fact that there situations not considered earlier. We choose

$`Z(x,t)`$ $`=`$ $`(M(x,t)+Q(x,t)^2)/Q(x,t),`$
$`M(x,t)`$ $`=`$ $`Q(x,t)^2/2,`$
$`R(x,t)`$ $`=`$ $`(2Q(x,t)Z_1(x,t)+Q(x,t)^3)/2.`$ (2.52)
By substituting all these expressions into the ansatz (2.50), we obtain the desired solution. The theorem is proved.

There is the following analogy with the theory of Korteweg-de Vries equtaions. In this context, relations (2.48) can be considered as an analog of Miura transformation.

Let us prove that all known solutions of the FhNS equation can be calculated in terms of the solutions of the a linear parabolic equation. We find the solution of the Riccary equation (2.48) and, choosing the upper sign, obtain

$$Q(x,t)=\frac{U(x,t)\mathrm{exp}(a_1x)}{C_0(aU(x,t)\mathrm{exp}(a_1x)dx}.$$
(2.53)
The compatibility condtions impose strong restrictions. Hence, for example, we have the pair of functions

$`Q(x,t)`$ $`=`$ $`{\displaystyle \frac{1+C_1\mathrm{exp}(2ax)}{1C_1\mathrm{exp}(2ax)}},`$
$`U(x,t)`$ $`=`$ $`C_1+\mathrm{exp}(2ax),a_1=a.`$ (2.54)
the substitution of which into (2.49) implies the following solution of Eq. (2.34):

$$u(x,t)=\frac{1+C_1\mathrm{exp}(2ax)}{1+C_1\mathrm{exp}(2ax)+\mathrm{exp}(ax3t\phi _3/2)}.$$
(2.55)
And the solution of (2.47)

$`U(x,t)`$ $`=`$ $`\left(1+2aC_2\mathrm{exp}(2ax)\right)\mathrm{exp}\left((a+a1)x+(aa1)(a+a1)t\phi _3/(2a^2)\right),`$
$`a_1`$ $`=`$ $`Ia\sqrt{2}.`$ (2.56)
gives solution of the equation (2.34)

$$u(x,t)=\frac{2aC_2+\mathrm{exp}(2ax)}{2aC_2\mathrm{exp}(2ax)+\mathrm{exp}(ax3t\phi _3/2)}.$$
(2.57)
These solutions have a specific characteristic feature, which can either appear or disappear. Such solutions were considered in , p. 64, , . In what follows, we consider Example 3 constructed by this method. In this example, two singularities simultaneously appear in the modified FhNS equation.

There is the following statement: to the WELL-KNOWN SOLUTION of Eq. (2.34)

$$u=\frac{1\mathrm{exp}(2ax)}{1+\mathrm{exp}(2ax)+\mathrm{exp}(ax3t\phi _3/2)},$$
(2.58)
there corresponds the solution of the LINEAR PARABOLIC EQUATION (2.46)

$`U(x,t)`$ $`=`$ $`(1+\mathrm{exp}(2ax))\mathrm{exp}((a+a1)x+(aa1)(a+a1)t\phi _3/(2a^2)),`$
$`a_1`$ $`=`$ $`Ia\sqrt{2}.`$ (2.59)
It turns out that the complex solutions of the linear parabolic equation can be recalculated into the solutions of the FhNS equation and vice versa. Therefore, the complex solutions of these equations must be studied.

Let us consider the FitzHygh–Naguma–Semenov (FhNS) equation (2.34) with variable coefficients.

By Theorem 2.1, for the first mode and separately for the third mode, there exist the following solutions. We have the following theorem.

We have the following theorem.

###### Theorem 2.3

Let $`k_0=k_0(x,t)`$, $`\phi _1=\phi _1(x,t)`$, $`\phi _3=\phi _3(x,t)`$ be smooth functions. Suppose that the function $`M(x,t)`$ is a solution of the equation

$$M_tk_0M_{xx}\phi _1M+\phi _3M^3=0,$$
(2.60)
and the following compatibility condition imposed on its coefficients is satisfied:

$`\sqrt{2}\left({\displaystyle \frac{M\phi _3k_{0}^{}{}_{t}{}^{}+k_0(2\phi _3M_t+\phi _{3}^{}{}_{t}{}^{}M)}{k_0\sqrt{m_2}}𝑑x}\right)\sqrt{m_2}+6M^2\phi _3\sqrt{m_2}4C_0^{}(t)\sqrt{m_2}`$
$`+6\sqrt{2}M_xk_0\phi _3\sqrt{2}Mk_{0}^{}{}_{x}{}^{}\phi _3+\sqrt{2}Mk_0\phi _{3}^{}{}_{x}{}^{}=0,`$ (2.61)
where $`m_2=k_0(x,t)\phi _3(x,t)`$ Then the function

$$u(x,t)=\frac{\mathrm{exp}(f)M(x,t)}{1\mathrm{exp}(f)},f=C_0(t)+\frac{M(x,t)\sqrt{\phi _3(x,t)}}{\sqrt{2k_0(x,t)}}𝑑x.$$
(2.62)
is a solution of the FhNS equation (2.34).

Proof. The proof follows from Theorems 1.2 and 2.1 and can be carried out by a direct substitution of the ansatz of the first mode, taking into account the fact that the coefficients in the equation are variable and $`M=0`$. One compatibility condition is identically satisfied. One compatibility condition remains. Hence the set of solutions is not empty. The following solution illustrates this theorem.

For example, we set $`k_0=1`$, $`\phi _1=\phi _3`$. We also assume, that we want to dress the solution $`Q=1`$. The compatibility conditon (2.62) turns into the equation for the function $`\phi _1`$:

$`\phi _{1}^{}{}_{t}{}^{}3\sqrt{2\phi _1}\phi _{1}^{}{}_{x}{}^{}+(\phi _{1}^{}{}_{x}{}^{})^2/\phi _1+\phi _{1}^{}{}_{xx}{}^{}=0.`$ (2.63)
We choose one of its solutions, for example, $`\phi _1=C_2\mathrm{tanh}^2(x/2\sqrt{3C_2\sqrt{2}})/3\sqrt{2}`$.

The solution of the FhNS equation

$$u_tu_{xx}\phi _1(uu^3)=0,$$
(2.64)
has the form

$$u(x,t)=\frac{f_1}{1f_1},f_1=\mathrm{exp}\left(\frac{C_2t}{2\sqrt{2}}\right)\mathrm{cosh}^{1/3}\left(\frac{\sqrt{3C_2}x}{2^{3/4}}\right).$$
(2.65)
This solution is real. If the constant is $`C_2>0`$, then this solution describes wave interaction . At the initial time moment, the solution has a singularity. Then the solution becomes smooth. For negative values of this constant $`C_2<0`$, this solution describes a periodic linear structure .

Thus one can construct a solution of the equation with variable roots of the source-sink function. For example, if $`\phi _1`$ is given to be a periodic function, then $`\phi _3`$ is constructed from solutions of Mathieu and Hill type equations.

###### Theorem 2.4

Let $`k_0=k_0(x,t)`$, $`\phi _1=\phi _1(x,t)`$, $`\phi _3=\phi _3(x,t)`$ be smooth functions. Suppose that one of the compatibility conditions on the coefficients is satisfied, namely, the first compatibility equation

$$Q_x(x,t)=\frac{2\sqrt{(}2\phi _3k_0)(UQ)^24k_0UQ(\phi _1UU_x)}{4k_0U^2},$$
(2.66)
or the second compatibility condition

$$Q_x(x,t)=\frac{2\sqrt{k_0}UQ2UQ^2\phi _3+2\sqrt{k_0}QU_x}{\sqrt{2k_0}U}$$
(2.67)
Suppose that the function $`Q(x,t)`$ is a solution of the FhNS equation

$$Q_tk_0Q_{xx}\phi _1Q+\phi _3Q^3=0,$$
(2.68)
and the function $`U(x,t)`$ is a solution of the linear parabolic equation

$$U_t+2\phi _1k_0U_xk_0U_{xx}U\left(\phi _{1}^{}{}_{t}{}^{}𝑑x+\phi _1+k_0\phi _1^2k_0\phi _{1}^{}{}_{x}{}^{}\right)=0.$$
(2.69)
Then the function

$$u(x,t)=\frac{Q(x,t)U(x,t)}{Q(x,t)\mathrm{exp}(\phi _1𝑑x)U(x,t)},$$
(2.70)
is a solution of the FhNS equation (2.34).

Proof. This statement can be proved by the substitution of the ansatz of the solution of the third mode (2.49), just as in the proof of Theorem 2.2. This theorem is a partial case in the analysis of the one of the two branches, which are determined by the choice of the compatibility condition.

The method has the following advantages:

a) the algorithm is iterative;

b) at each iteration step we solve one ordinary differential equation of the first order or one algebraic equation (cf. , where one has to solve a system of ordinary differential equations);

c) the method allows one to obtain new solutions; it is necessary only to know which mode has nontrivial solutions; to this end, there are some considerations the discussion of which is beyond the framework of this paper.

Note that formulas (2.41), (2.42) can be also obtained in a different way, while formulas (2.49), (2.63), (2.72) are constructed by this method.

EXAMPLE 2. Let us consider a modified FitzHygh–Nagumo–Semenov equation, which is a special case of the equation (1.1). In this case there also exist a lotof solutions, but we write only one solution correspondign to the second mode, since it differs from all other solutions. The change (1.2) reduces our equation to an equation the fourth order of homogeneity.

Suppose that Eq. (1.1) has the form

$$(a_1h_1u)u_t(a_1h_1u)u_{xx}h_1(a_1+h_1u)u_x2h_1(u_x)^2a_1(a_1^2+h_1^2)u(1u^2)/2=0,$$
(2.71)
The solution has the form

$$u(x,t)=\frac{1\mathrm{exp}(a_1x)}{1+\mathrm{exp}(a_1x)+\mathrm{exp}(((h_1^23a_1^2)t+2(h_1+a_1)x)/4)}$$
(2.72)
and describes the development of the second wave from small perturbations.

This solution can be used for modeling phenomena related to phase transitions . One of the elementary equations for which this effect is preserved for $`a_1=0`$ has the form

$$u_tu_xx+h_1u_x+2(u_x)^2/u=0,$$
(2.73)
EXAMPLE 3. We consider the modified FhNS equation (1.3):

$$u_t+2(h_1)^2u_{xx}/\phi _3H_1(1+4u)u_x\phi _3u(1u^2)=0,$$
(2.74)
The solution has the form

$$u(x,t)=\frac{1\mathrm{exp}((h_1t+x)\phi _3/h_1)}{1\mathrm{exp}(x\phi _3/(2h_1))+\mathrm{exp}((h_1t+x)\phi _3/h_1)},$$
(2.75)
This exact solution shows how a solution of the problem with a smooth initial condition of double singularity is developed during a finite time interval. A discontinuity arises at $`t=1.4`$. This discontinuity is of a special structure such that the left-hand discontinuity does not move, while the right-hand discontinuity moves to the right.

Note that if the sign of the parameter $`\phi _3`$ is changed (i.e., we have $`\phi _3=1`$), then Eq. (1.1) becomes an inverse parabolic equation. However, the solution (1.1) exists and is bounded. Small perturbations of the initial condition relax in time.

EXAMPLE 4. We consider the solution correspondign to the second mode for the equation the fourth order of homogeneity. Suppose that the modified FhNS equation (1.1) has the form

$$(1u)u_t(1u)u_{xx}+(1/a_1+u/a_12a_1u\phi _3)u_x2(u_x)^2\phi _3u(1u^2)=0.$$
(2.76)
This equation has a solution of the form

$$u[x,t]=\frac{1\mathrm{exp}(x/a1)}{1+\mathrm{exp}(x/a1)+\mathrm{exp}(x/a1t\phi _3)}.$$
(2.77)
The solution shows how the structure develops from an initial state whose range is $`[0.5,1]`$ to a traveling wave $`u[0,1]`$. Note that it is possible to construct many solutions of such type, i.e., solutions that describe how a certain wave evolves from some intial state.

EXAMPLE 5. We consider the modified Fisher–Kolmogorov–Petrovskii–Peskunov equation

$$u_tk_0u_{xx}(2k_0\phi _2/h_2+h_2u)u_{xx}\phi _1u+\phi _2u^2=0.$$
(2.78)
We consider an important example of solutions to the Fisher–Kolmogorov–Petrovskii–Piskunov modified equation and show that its solutions can be expressed in terms of higher transcendental functions such as Bessel functions or hypergeometric functions. We also show how they are related to the Bäcklund transformation for the Bürgers equation.

###### Theorem 2.5

Suppose $`b_1=0`$, $`\phi _3=0`$,

$$h_1=2k_0\phi _2/h_2$$
(2.79)
in equation (1.3). The function $`U(x,t)`$ is a solution of the linear parabolic equation

$$U_t=r(x,t)U+k_0U_{xx},$$
(2.80)
with the potential

$$r(x,t)=\mathrm{exp}(2x/h_2+C_1(t)).$$
(2.81)
The function $`U(x,t)`$ has the form

$$U(x,t)=v(z(x,t)),Z(x,t)=\mathrm{exp}(x/(2h_2)+C_1(t)/2))$$
(2.82)
and the function $`v(z)`$ is a solution of the equation

$$2h_2^2z\phi _2v(z)+h_2^2\phi _1v^{}(z)+2k_0\phi _2v^{}(z)+2k_0z\phi _2v^{\prime \prime }(z)$$
(2.83)
Then the function $`u(x,t)`$ is a solution of Eq. (2.78) and has the form

$$u(x,t)=(2k_0/h_2)_x(\mathrm{ln}(U(x,t))),$$
(2.84)
Proof. We substitute (2.84) into the original equation (2.78) and obtain a nonlinear homogeneous equation from which we exclude the derivatives $`U_t`$, $`U_{xt}`$ with the help of Eq. (2.80). Then, for the parameters specified in the assumptions of the theorem by the change (2.83), we obtain ordinary defferential equations for higher transcendental functions. The proof of Theorem 2.5 is complete. $`\mathrm{}`$

It should be noted that the change (2.84) coincides in form with the Cole–Hopf change of variables for the Bürgers equation.

The author is grateful to V. G. Danilov and S. Yu. Dobrohotov for constant attention to his work and useful discussions and to V. P. Maslov, and A. D. Polynin for constructive advice.


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# Comparative Chemistry of Diffuse Clouds II: CN, HCN, HNC, CH3CN & N2H+

## 1 Introduction.

The presence of interstellar CN was shown by McKellar (1940) who identified the R(0) line at $`\lambda `$387.4nm (see also Swings & Rosenfeld (1937) and Adams (1941)). The R(1) line being undetectable, and taking as an upper limit on its equivalent width a value one-third that of R(0), he calculated the “effective” or “rotational” temperature of interstellar space to be no more than 2.7 K. McKellar qualified his conclusion, inquiring “if, indeed, the concept of such a temperature in a region with so low a density of both matter and radiation has any meaning.” Of course it does have meaning, and CN has long been used to derive the temperature of the cosmic microwave background at a wavelength of 2.7mm, the contribution to the rotational excitation from collisions often being negligibly small (Thaddeus, 1972; Crane et al., 1989; Black & van Dishoeck, 1991). Both CH and CH+, the other early molecules, are too light for their rotational structure to have been used in this way.

Sixty years after the discovery of the CN lines we are still finding out what their presence in diffuse clouds actually means. Although it has not generally been regarded as much of a challenge to come up with rather conventional chemical schemes which account for the observed abundances of CN (Federman et al., 1994; Van Dishoeck & Black, 1986), these models beg larger questions which are the subject of this series of papers. Earlier we showed that the abundances of the closely related species HCN and HNC are much larger than expected in diffuse clouds seen toward compact extragalactic mm-wave continuum sources at galactic latitudes of 2-15o (Lucas & Liszt, 1993, 1994). Given that the proximate source of CN in conventional models is the breakup of HCN and/or HNC and their immediate antecedents, it follows that our understanding of CN is at best accidental (see also Crawford & Williams (1997)).

In the first paper in this new series (Lucas & Liszt, 2000), we discussed the comparative chemistry of species in the CnHm-family of hydrocarbons whose most abundant known members are CH, C2H and C3H2. Here, we present results for the cyanogen- or CN-bearing molecules CN, HCN and HNC, based on a modest survey of $`\lambda 3`$mm absorption from clouds occulting a subset of our wider sample of compact extragalactic continuum sources (Lucas & Liszt, 1996). We find a remarkably tight linear correlation among all three, such that it is hardly necessary to observe the other two once the abundance of any one has been found: a similarly tight linear relation exists between OH and HCO+ (Liszt & Lucas, 1994; Lucas & Liszt, 1996; Liszt & Lucas, 2000) and a somewhat looser but still quite good one exists between C2H and C3H2. By contrast, intercomparisons among the different families (HCN or C2H $`vs.`$ HCO+; HCN $`vs.`$ C2H) reveal a high degree of correlation, but sharply non-linear behaviour. In general, all species except OH increase in abundance more rapidly than linearly with increasing N(HCO+) at N(HCO+) $`10^{12}\mathrm{cm}^2`$ . Analogous non-linear behaviour in the optical regime is seen in plotting N(CN) $`vs.`$ N(C2) (power-law slope around 1.4, see below) or either of these against N$`(\mathrm{H}_2)`$, but tight linear relationships between species are quite lacking in optical absorption studies.

The manner of taking and analyzing the present observations is discussed in Section 2. The main results of this work are discussed in Section 3, and a discussion of what is known of the chemistry of simple CN-bearing molecules is presented in Section 4.

## 2 Observations.

The mm-wave absorption data discussed here were taken at the Plateau de Bure Interferomer during the years 1993 - 1997. Table 1 shows the list of background sources observed, along with their galactic coordinates and the rms noise levels achieved in the line/continuum ratio, which is also the rms error in optical depth in the optically thin limit. The mode of observation was identical to that described previously by Lucas & Liszt (2000) and in our other prior work. All data used for the present survey were taken with a channel separation of 78 kHz and a channel width (FWHM) of 140 kHz, corresponding to velocity intervals which are given just below for the various individual species.

We have never succeeded in detecting emission from HCN, the most optically thick of the species discussed here, from the clouds occulting any of our extragalactic continuum sources. Thus we derive all column densities in the limit of no collisional excitation, i.e. $`T_{\mathrm{exc}}`$ = $`T_{\mathrm{cmb}}`$. Very weak HCN emission (0.02 K) was seen some 30 South of $`\zeta `$ Oph at the local HCO+ emission peak (0.1 K; Liszt (1997)) but CN emission is absent at very low levels toward the star itself (Crane et al., 1989; Kopp et al., 1996).

### 2.1 CN

We observed the (1/2-1/2), (3/2-5/2) and (1/2-3/2) hyperfine components of the N=0-1 absorption line near 113.490GHz, which have LTE line strengths in the proportion 988:3333:1235. As shown in Figure 1, the lines are split by 3-11 MHz and so are separately resolved. Assuming equilibrium with the cosmic microwave background, the CN column density is related to the integrated optical depth of the N=0-1 absorption as N(CN) $`=2.27\times 10^{13}\mathrm{cm}^2\tau _{01}𝑑v`$.

Results of Gaussian fitting are given in Table A1 (see Section 2.6). The channel separation and FWHM of the data are 0.206 km s-1 and 0.370 km s-1, respectively.

### 2.2 HCN

We observed the F=1-1, F=2-1 and F=1-0 hyperfine components of the N=0-1 transition near 88.63GHz, with LTE line strengths in the ratio 1:5:3. As shown in Figure 1, the lines are split by 1-3 MHz and are separately resolved, though often confused due to blending with other features. Assuming equilibrium with the cosmic microwave background, the HCN column density is related to the optical depth of the N=0-1 absorption integrated over all components as N(HCN) $`=1.89\times 10^{12}\mathrm{cm}^2\tau _{01}𝑑v`$. Results of Gaussian fitting are given in Table A2 (see Section 2.6). The channel separation and FWHM of the data are 0.264 km s-1 and 0.474 km s-1, respectively.

### 2.3 HNC

We observed the F=1-1, F=2-1 and F=1-0 hyperfine components of the N=0-1 transition near 90.66GHz, with LTE linestrengths in the ratio 1:5:3. As shown in Figure 1, the hyperfine splittings of the J=0-1 absorption of HNC are smaller than typical linewidths, and so are observed as a broadening of the observed components rather than as separate kinematic features. Assuming equilibrium with the cosmic microwave background, the HNC column density is related to the integrated optical depth of the J=0-1 absorption as N(HNC) $`=1.76\times 10^{12}\mathrm{cm}^2\tau _{01}𝑑v`$. Results of Gaussian fitting are given in Table A3 (see Section 2.6). The channel separation and FWHM of the data are 0.256 km s-1 and 0.463 km s-1, respectively.

### 2.4 CH3CN

We observed a band centered at 91.9871 GHz, $`i.e.`$ at the frequency of the strongest hyperfine component of the J=5-4 transition in the K=0 ladder, toward B0415+379, B0528+134, B1730-130, and B2251+158. This band sampled as well the same J-transition in the K=1 ladder, whose intensity should be smaller by a factor 0.076 in equilibrium with the cosmic microwave background; yet higher K-ladder transitions falling in the band should have truly negligible population. Although the spectrum is mildly suggestive of a detection toward 3C111, it is not statistically significant above the $`2\sigma `$ level. The spectrum toward 3C111 had an rms noise in the line/continuum ratio of 0.004, leading to a 2-sigma upper limit on the integrated optical depth of $`\tau 𝑑v<0.014`$ km s-1, or N(CH3CN) $`<2.4\times 10^{11}\mathrm{cm}^2`$, N(CH3CN)/N(HCO+) $`<`$ 0.021. These are the results quoted in Table 3 and discussed below. CH3CN is underabundant compared to the composition of dark clouds like TMC-1 and L134N.

### 2.5 N2H+

We observed a band around 93.1738 GHz, the location of the J=1-0 transitions of N2H+, toward B0355+508, B0415+379, B0528+134, B1730-130, and B2251+158, but detected no absorption. The best limit is for the gas occulting B0415+379, and this is the datum used in Table 3 and discussed below. The spectrum toward 3C111 had an rms noise in the line/continuum ratio of 0.0025, leading to a 2-sigma upper limit on the integrated optical depth of $`\tau 𝑑v<0.016`$ km s-1, or N(N2H+) $`<2.6\times 10^{10}\mathrm{cm}^2`$, N(N2H+)/N(HCO+) $`<0.0024`$. By comparison, the other lines of sight did not furnish useful data except to confirm that we had not somehow missed large quantities of N2H+ anywhere. N2H+ is much more strongly underabundant in diffuse gas than CH3CN.

### 2.6 Component fitting

In Tables A1-3 of the Appendix we give the results of Gaussian decomposition, done simultaneously to whichever subset of the hyperfine structure was actually observed (all for HNC and HCN, and part for CN, see Figure 1). The optical depth quoted at line center pertains to the strongest hyperfine component but was derived from a simultaneous fit to multiple kinematic and hyperfine components, the latter assumed to appear in the LTE ratio, as discussed in Section 2.4. The integrated optical depth tabulated is the sum over all hyperfine components for a given kinematic feature, either derived directly (for HCN and HNC) or (for CN) by a simple scaling to account for that (smaller) fraction of the spectrum which was not actually observed.

We always fit the observed line profile (line/continuum ratio) using saturated gaussians of the form $`\mathrm{exp}(\tau _0\times \mathrm{exp}((vv_0)^2/2\sigma _{v}^{}{}_{}{}^{2}))`$, simultaneously considering all the accessible hyperfine structure. Thus the optical depths and linewidths which are presented in the Tables are consistent with each other and with the degree of saturation which occurs in the data. The fitted linewidths are not artifically increased by line saturation.

As seen in Figure 2, the intrinsic kinematic structure along a given line of sight can introduce substantial confusion as to which clouds contain which molecule. In each case, we used our HCO+ profile as a template and did a simultaneous fit, species-by-species, to all of the hyperfine and kinematic components which were observed. The FWHM linewidths in Tables A1-A3 of the Appendix all include a contribution from the intrinsic broadening in the spectrometer which, as noted above, has channel half-power widths somewhat wider than the channel separation itself. Where the linewidths are plotted (Figure 5), this contribution has been subtracted in quadrature from the tabulated value.

## 3 Systematics

### 3.1 Abundance and linewidth comparisons within the family of small CN-bearing molecules

Figure 3 at top shows that N(CN), N(HCN), and N(HNC) have the kind of tight linear relationship previously demonstrated for OH and HCO+ (Liszt & Lucas, 1996; Lucas & Liszt, 1996). The slopes of both regression lines in Figure 3 are indistinguishable from unity ($`1.03\pm 0.08`$ in either case) with mean ratios (weighted by the variance) of $``$N(HNC)/N(HCN)$`=0.21\pm 0.05`$ and $``$N(CN)/N(HCN)$`=6.8\pm 1.0`$.

Figure 5 shows a comparison of the FWHM linewidths derived from profile decomposition (see the Tables in the Appendix), corrected for finite resolution. The HCN, HNC, and HCO+ linewidth intercomparisons show no deviation from the condition of equal linewidths, but CN seems to have a slight tendency to be narrower than HCN and a more marked tendency to be narrower than HCO+. CN is typically observed to be narrower than CH in optical spectra, where the difference is usually interpreted as suggesting a spatial segregration of CN into the denser core regions of a cloud (Crawford, 1995).

In this work, CN was observed with a somewhat better spectral resolution, owing to its higher rest frequency and the fixed nature of the channel widths in the autocorrelator. Although it seems unlikely that resolution dominates the linewidth comparison in Figure 5 – HCN and HCO+ are observed at nearly the same frequency and resolution – this should certainly be checked. If CN really is narrower than HCN, this is extremely interesting because the tight correlation of their column densities militates against any really large-scale spatial segregation, such as isolating CN in a denser core. Instead, there would have to be a mix of species and conditions throughout the gas such that line of sight average column densities vary in fixed proportion, and the CN linewidth would be smaller throughout the cloud, or at least down to scales much finer than we can discern directly.

### 3.2 Abundance comparisons across families of species observed at mm-wavelengths

Figure 3 in the middle shows the variation of HCN with respect to HCO+ and C2H. As we have shown before, comparisons across chemical families reveal sharply non-linear behaviour; although the variations of HCN with HCO+ and C2H may appear similar, the C2H-HCO+ relationship is also non-linear, with the same increase in N(C2H) at N$`(\text{HCO}\text{+})1\times 10^{12}\mathrm{cm}^2`$ (Paper I). C2H however has something of a plateau at small N(HCO+), where the mean of N(C2H)/N(HCO+) is high. The N(HCN)/N(C2H) ratio varies from (roughly) 0.08 to 0.3, increasing with column density.

Taken as unweighted averages over all observed features, we find $`<`$N(HCN)/N(HCO+)$`>=1.47\pm 0.86`$, $`<`$N(HCN)/N(C2H)$`>=0.155\pm 0.094`$; in Paper I we quoted $`<`$N(C2H)/N(HCO+)$`>=14.5\pm 6.7`$ for the unweighted mean. When the averages are weighted by the fractional errors of the column density, favoring higher column density sources, we have $`<`$N(HCN)/N(HCO+)$`>=1.97\pm 0.43`$ and $`<`$N(HCN)/N(C2H)$`>=0.23\pm 0.57`$.

### 3.3 Insights from optical absorption data

Figure 3 at bottom shows optically-determined CN, C2, and $`\mathrm{H}_2`$ column densities from the updated version of the summary compilation of Federman et al. (1994). It is lamentable that there are so few lines of sight for which both molecular hydrogen and CN have been observed. In any case, although the relationships in Figure 3 at the bottom are non-linear (power-law slopes of 1.42 and 1.31 reading from left to right) they provide a useful calibration of absolute abundances. For instance, at N(CN) = $`2.1\times 10^{12}\mathrm{cm}^2`$, the fitted curve gives N($`\mathrm{H}_2`$) = $`4.2\times 10^{20}\mathrm{cm}^2`$. This can be transferred to the fits at the top of Figure 3 to infer N(HCN) = $`3.1\times 10^{11}\mathrm{cm}^2`$ and then to Figure 3 at the middle for (very) approximately N(HCO+) $`7.2\times 10^{11}\mathrm{cm}^2`$, or N(HCO+)/N($`\mathrm{H}_2`$) $`1.7\times 10^9`$. This is near the single value N(HCO+)/N($`\mathrm{H}_2`$) = $`2\times 10^9`$ that we have used for a variety of purposes recently (Liszt & Lucas, 2000), most especially to show that the observed amounts of CO in diffuse clouds follow directly from recombination of this amount of HCO+ with electrons in standard cloud models.

As inferred from this train of thought, the HCO+ abundance is found to be slightly smaller ($`1.5\times 10^9`$) at N($`\mathrm{H}_2)=1.0\times 10^{21}\mathrm{cm}^2`$. But from this comparison we gain the new insight that the abrupt increases in the abundances of so many species at N$`(\text{HCO}\text{+})10^{12}\mathrm{cm}^2`$ are actually occurring at N($`\mathrm{H}_2)0.30.5\times 10^{21}\mathrm{cm}^2`$.

### 3.4 Linear and non-linear functional dependencies

It is important to understand which parameters and physical conditions are varying in Figure 3. For instance, at the top there is a factor 50 variation in N(CN), N(HCN) and N(HNC), but much of this – perhaps a factor 20-30 – occurs over a much smaller range in N(HCO+): a factor 5-6. Thus the total column density of gas probably changes far less than might otherwise be guessed. At the bottom, N(CN) seems to vary by a factor ten or so at N($`\mathrm{H}_2`$) $`3\times 10^{20}\mathrm{cm}^2`$, and less rapidly, but still faster than linearly, at higher N($`\mathrm{H}_2`$).

To tie everything together somewhat more tightly, we show in Figure 4 the variation of N(CN) and N($`\mathrm{H}_2`$) with total gas column density N(H) = N(H I) + 2N($`\mathrm{H}_2`$) (Savage et al., 1977) for those sightlines in Figure 3 where CN has been measured at optical wavelengths. The point is that, while N(CN) and N($`\mathrm{H}_2`$) are strongly coupled at higher N($`\mathrm{H}_2`$) in Figure 3, neither N(CN) nor N($`\mathrm{H}_2`$) increases with N(H) or extinction along those same directions. Both N(CN) and N($`\mathrm{H}_2`$) vary substantially over a very small range of N(H). 111In general the effect of extraneous atomic hydrogen is a rightward displacement in Figure 4, causing the data to appear consistent with $`\mathrm{H}_2`$-formation at too low a local number density. The anomalous point at large N(H) in Figure 4 for the sightline to $`\rho `$ Oph A is an extreme example. No other sightline in the Copernicus sample had such a large N(H), for any value of N($`\mathrm{H}_2`$).

The theoretical $`\mathrm{H}_2`$-formation curves superposed in Figure 4 at left were constructed by sampling along the mean line of sight through a uniformly-illuminated spherical gas parcel of constant total density n(H) (Liszt & Lucas, 2000). They suggest that those sightlines with larger N($`\mathrm{H}_2`$) – and so with larger N(CN) – sample denser and perhaps darker gas with a higher local fraction of molecular hydrogen. This in turn leads to the idea that there is some cascade of effects which accounts for the very strong changes in column density among the trace species, largely driven by the final approach of the molecular hydrogen fraction toward unity, but perhaps having as its cause the change in some other agent.

In the most classical terms, an increase in local number density would drive the molecular hydrogen fraction higher at relatively constant extinction, increasing the rate of formation of ancestor molecular ions like $`\mathrm{H}_2`$CN+, which would then recombine faster at higher density and (if the thermal pressure is even approximately conserved, as is almost the case in our models) somewhat lower temperature. It helps considerably that both the electron fraction and the abundance of the basic building material (C+) are dominated by the C+/C0 ratio, which remains high in diffuse gas at the moderate density and extinction relevant here. Earlier we showed that even a constant abundance (relative to $`\mathrm{H}_2`$) of an ancestor ion could give rise to faster-than-linear increases in the abundance of a product molecule via electron-ion recombination in a gas of constant density, specifically N(CO) varying as N$`(\mathrm{H}_2)^2`$ when arising from gas with constant N(HCO+)/N$`(\mathrm{H}_2`$) (Liszt & Lucas, 2000).

A related matter of interest is the extent to which the molecules we sample may exist in environments which are still partly atomic. In Figure 6 we show the variation of the $`\mathrm{H}_2`$-fraction with position inside spherical gas parcels of density n(H) = 16 and 128 H-nuclei $`\mathrm{cm}^3`$, for five values of the column density across the central line of sight through the parcel. At the panel to the right in this Figure, where the number density is large enough that weak HCO+ emission might be seen, even low-extinction gas is mostly molecular: the lowest column density considered here is about that of a standard “Spitzer” HI cloud. But at low density, and when small molecular column densities are sampled, the ambient environment of hydrogen gas is likely to be largely atomic.

## 4 Chemistry of simple nitrogen-bearing molecules

### 4.1 CN, HCN and HNC

Shown in Table 2 is a comparison across diffuse and dark gas of the chemical abundances of the hydrocarbons studied in Paper I and species relevant to the present work (NH, CN, HCN, HNC). The tabulation indicates that the abundances of HCN and HNC increase very substantially in moving from diffuse to dark gas while CN increases in abundance only slightly.

As represented in the last column of Table 2 by the work of Van Dishoeck & Black (1986), models of moderate-density quiescent gas-phase diffuse cloud chemistry provide approximately the correct amounts of NH and CN, but fail to provide the inferred amount of HCN (which we derive by scaling N(CN)) by a factor of 10. However even this degree of agreement may be illusory. The Van Dishoeck & Black (1986) results were computed for a very much larger carbon abundance than is now accepted (see the entry under $`\zeta `$ Oph from Cardelli et al. (1993)) and the claim is made (see the discussion of Crawford & Williams (1997)) that ordinary gas phase diffuse cloud chemistry fails by a factor of about 40 to reproduce the observed N(NH) and N(CN). If so, the predicted amounts of most molecules would be much smaller than are indicated in the last column of Table 2.

It seems generally accepted that formation of the CN-bearing molecules proceeds $`via`$ the reaction C$`{}_{}{}^{+}+\mathrm{NH}\mathrm{CN}^++`$H, followed by a series of hydrogen abstractions beginning with CN$`{}_{}{}^{+}+\mathrm{H}_2\mathrm{HCN}^++`$ H, and terminating (forming one of the observed species CN, HCN, HNC) after one of many possible electron-ion recombinations such as $`\mathrm{H}_2`$CN$`{}_{}{}^{+}+\mathrm{e}\mathrm{HCN}+`$H (Van Dishoeck & Black, 1986; Federman et al., 1994; Crawford & Williams, 1997). However, the particulars are very much open to question since the needed gas-phase reactions do not proceed nearly fast enough under typical quiescent conditions in diffuse clouds.

Quiescent gas phase chemical models make NH by cosmic ray ionization of N, followed by hydrogen abstraction in the reaction of N+ and $`\mathrm{H}_2`$. This and all other gas phase NH formation reactions are discounted by Crawford & Williams (1997) in favor of formation on grains; they also argue that the linewidth of NH, which is typical, rules out any formation mechanism involving translationally hot particles in the gas phase, like ions (N+) accelerated in MHD shocks.

Whether the NH may be formed in gas-phase reactions with fast-moving particles is of interest because a similar formation route is cited as the likely source of the unexplained high abundances of HCO+ ($`i.e.`$ C+ \+ OH $``$ CO+ \+ H) and CH+ ($`via`$ C+ \+ $`\mathrm{H}_2`$) (Falgarone et al., 1995; Hogerheijde et al., 1995; Federman et al., 1996; Joulain et al., 1998). One of these pathways also seems to be needed to explain the observed amounts of CO, as well, and they would help greatly in alleviating the apparent shortfall of the CN-bearing species (through the reaction of C+ and NH). If such processes do not contribute to alleviating the apparent shortfall of NH, it seems possible that they might not be needed at all to form the CN-bearing species. Alternately, if they do function, they would represent strong sources and sinks of NH not considered by Crawford & Williams (1997).

Whether the linewidth of NH rules out an enhanced gas phase formation rate really awaits a calculation of reaction dynamics and subsequent molecular thermalization in the ambient medium (Bucher & Glinski, 1999). For molecules which are formed hot – which includes both those which are formed from translationally hot reactants and those which carry away significant amounts of energy from exoergic chemical reactions – only those which react quickly will disappear before reaching kinetic equilibrium and so exhibit a broader linewidth. CH+, for instance, reacts rapidly with both $`\mathrm{H}_2`$ and electrons, while NH does not.

### 4.2 CH3CN

Table 2 shows that CH3CN is underabundant in diffuse gas relative to TMC-1, by at least a factor of twenty, or somewhat less, 6, if its abundance relative to HCO+ is considered. Given that easily-observed species like C3H2 and HCN are underabundant by a similar factor, the limit on CH3CN is not that strong and should probably not be considered too surprising. Both of the chemical sequences leading to CH3CN in dark gas (Nejad & Wagenblast, 1999), (C, C$`\mathrm{H}_2`$, HCN, H4C2N+, CH3CN) at early times or (C3H, C3N, HC3N+, HCN, H4C2N+, CH3CN) later on, run through HCN.

### 4.3 N2H+

N2H+ is one of several molecules (N2H+, NH3 , SO2) which are considered to form only very late in the chemical evolution of dark clouds (Nejad & Wagenblast, 1999). Observationally, NH3 at least is much more widespread than N2H+ and has been seen toward B0355+508 by Nash (1990) and toward 3C111 in our own unpublished work. However, unpublished VLA results communicated by A. Marscher indicate that reliable NH3 abundances await a survey with the upgraded K-band receivers at the VLA, and that neither Nash’s profiles nor ours (all taken with single dishes), are trustworthy. Unfortunately, NH3 column densities are also unavailable from Marscher’s VLA profiles, because only the (1,1) lines were observed.

In molecular clouds, N2H+ forms in appreciable quantities at late times $`via`$ H3+$`+`$ N2 $``$ N2H+ $`+`$ $`\mathrm{H}_2`$, i.e. when the ionization fraction has dropped and when CO has been depleted from the gas phase, encouraging a higher abundance of H3+. It is known for being detectable only in darker regions.

Toward 3C111, N(N2H+)/N(HCO+) $`<0.002`$ (2-sigma) as compared to 0.25 in TMC-1. This is by far the best limit we have ever been able to set on an undetected species. 3C111 is actually seen through a small hole (region of lower than average extinction) in an outlying cloud in the Taurus cloud complex; a discussion of the viewing geometry is given by Lucas & Liszt (1998).

## 5 Summary

We showed that the abundances of CN, HCN and HNC vary in fixed proportion 6.8:1:0.21 respectively. This is rather different from conditions in TMC-1, where their abundances are in the ratio 6.8:4.5:4.5. The fractional abundances of HCN and HNC increase substantially between diffuse and dark gas but that of CN increases only by factor of a few. Nonetheless, HCN and HNC are much more abundant in diffuse gas than predicted by models of quiescent low-density gas-phase chemistry, with an intrinsic variation over the approximate range N(HCN)/N($`\mathrm{H}_2)1.45.0\times 10^9`$ for N(HCO+)/N($`\mathrm{H}_2)=2.0\times 10^9`$. The latter value is confirmed by comparing our measurements of CN, HCN, and HCO+ with the lamentably slight body of optical absorption data for lines of sight where both CN and $`\mathrm{H}_2`$ have been observed.

The linewidths of HCN, HNC, and HCO+ are sensibly identical, but CN shows a slight tendency to be narrower than HCN and a more pronounced tendency to be narrower than HCO+. If the CN and HCN are strongly coupled in column density but different kinematically, this argues against a spatial segregation of CN into denser regions of a cloud (core), such as is usually invoked to explain the narrower width of CN absorption lines seen in the optical. Instead, clouds would have to be a mixture of kinematically (and chemically) distinct regions all occurring many times along any one line of sight.

In comparing the abundances of species in the CN family with the hydrocarbons discussed in Paper I (principally C2H and C3H2) or with HCO+, we find strong correlations but sharply non-linear variation. Both the CN-species and hydrocarbons show a strong increase in column density around N(HCO+) $`=10^{12}\mathrm{cm}^2`$ and all comparisons across families show relative abundance changes by factors of 3 or more. We note that C2H and C3H2 vary in fixed proportion with each other, though not as tightly as for HCO+ and OH, or as seen here.

For those sightlines where the column densities of CN, HI, and $`\mathrm{H}_2`$ have been measured optically, it appears that N(CN) turns on abruptly at N($`\mathrm{H}_2`$) $`3\times 10^{20}\mathrm{cm}^2`$ and increases faster than linearly at larger N($`\mathrm{H}_2`$), in a regime where the abundance of $`\mathrm{H}_2`$ itself varies widely over a very small range in N(H) (or extinction). Models of $`\mathrm{H}_2`$-formation suggest that the hydrogen is still largely atomic in the gas sampled on sightlines with the smallest column densities of the trace species like CN.

Earlier (Liszt & Lucas, 2000), we showed that N(CO) $`N(\mathrm{H}_2)^2`$ (approximately) when CO forms in a gas of constant density via recombination of HCO+, and the latter is held in fixed proportion ($`2\times 10^9`$) to $`\mathrm{H}_2`$. It follows that much stronger variations in trace abundances would occur if the higher molecular fractions along some sightlines – reflected in higher N($`\mathrm{H}_2`$) and N(HCO+) – arise because the density is higher locally; the sequence of molecular abundances revealed in our work must at least partly reflect the tendency for higher-density regions to harbor higher $`\mathrm{H}_2`$-fractions at a given total hydrogen column density. Of course this cannot be the whole story because the high abundances of the parent ions needed to reproduce our observations of their descendants cannot presently be explained by any agreed-upon physical scenario. Equally important is to understand why the abundances of some species relative to $`\mathrm{H}_2`$ vary so strongly ($`e.g.`$ C2H, C3H2, CN, HCN, and HCN) while those of CH, OH, and HCO+ seem nearly fixed.

We failed to find CH3CN and N2H+, meaning that they are at least 10 and 100 times less abundant (relative to HCO+) in diffuse clouds than in TMC-1. The limit on CH3CN may not be very significant, given that C3H2 and HCN are underabundant by about the same amount but are readily seen along many lines of sight and that, in dark gas, the formation paths to CH3CN pass through HCN at both early and late times. The limit for N2H+ may be more interesting but this species is notoriously hard to detect even in translucent clouds of much higher density than we study.

We discussed what is known of the chemistry of the CN-HCN-HCN species. Although there is a consensus that they form through some chemical pathway involving NH, most likely interacting with C+, the origin of NH and the mode of interaction are entirely uncertain.

The next papers in this series will discuss the sulfur-bearing species CS, SO, SO2, $`\mathrm{H}_2`$S and HCS+, and probably NH3 and $`\mathrm{H}_2`$CO.

###### Acknowledgements.

We thank Dr. John Black for a close reading of the manuscript. The National Radio Astronomy Observatory is operated by AUI, Inc. under a cooperative agreement with the US National Science Foundation. IRAM is operated by CNRS (France), the MPG (Germany) and the IGN (Spain). We owe the staff at IRAM (Grenoble) and the Plateau de Bure our thanks for their assistance in taking the data.

## Appendix A Products of Gaussian fitting


text:
# References

EDO-EP-40

April, 2001

Trapping of Nonabelian Gauge Fields on a Brane

Ichiro Oda 111 E-mail address: ioda@edogawa-u.ac.jp

Edogawa University, 474 Komaki, Nagareyama City, Chiba 270-0198, JAPAN

## Abstract

We show that as in abelian gauge fields, nonabelian gauge fields are also trapped on a brane in the Randall-Sundrum model by applying a new mechanism based on topological Higgs mechansim. It is pointed out that although almost massless gauge fields are localized on the brane by the new mechanism, exactly massless gauge fields are not localized. This fact does not yield any problem to abelian gauge fields, but may give some problem to nonabelian gauge fields since it is known that there is a discontinuity between massless and massive fields in the case of the nonabelian gauge theory as in the gravitational theory.

The brane world scenarios based on the gravity-localized models have opened new perspectives for elementary particle physics beyond the Standard Model . (For multi-brane models, see .) The problems in modern physics such as the cosmological constant, the supersymmetry breaking and the hierarchy problems might be solved within the framework of the brane world scenarios. These problems are hard since solving such problems seems to require a low energy mechanism, but not only the low energy physics in the standard framework of effective theory does not seem to offer a solution but also it is very difficult to change the low energy framework in a sensible manner, given all of the familiar successes of the Standard Model.

The key ingredient in the gravity-localized models with noncompact extra spaces is how to localize all the matter and gauge fields in addition to the graviton on a brane. For instance, if charged particles cannot be sharply localized on the brane, we are in contradiction with the well-established experiment of the charge conservation law in our world. Indeed, concerning the localization of various bulk fields on the brane (or the domain wall), there have been a lot of works thus far within the framework of local field theory -. (For a review, see .) In particular, the non-localization of bulk gauge fields on a brane was one of fatal drawbacks associated with the gravity-localized models since there certainly exists a massless ’$`\mathrm{𝑝ℎ𝑜𝑡𝑜𝑛}`$’ in our world.

Recently, we have proposed a new localization mechanism for abelian gauge fields . The idea of constructing such a new mechanism has stemmed from an attempt of making a gauge field’s analog of the localization mechanism of bulk fermions. The aim of the present paper is to apply the new mechanism to nonabelian gauge fields, i.e., ’$`\mathrm{𝑔𝑙𝑢𝑜𝑛𝑠}`$’. We will see that the nonabelian gauge fields are also localized on the brane by the same mechanism, but there is a subtle issue associated with a discrete difference between the zero-mass theories and the non-zero mass theories.

Let us first explain the model setup. We consider the following $`AdS_5`$ metric:

$`ds^2`$ $`=`$ $`g_{MN}dx^Mdx^N`$ (1)
$`=`$ $`\text{e}^{A(r)}\eta _{\mu \nu }dx^\mu dx^\nu +dr^2,`$
where $`M,N,\mathrm{}`$ are five-dimensional space-time indices and $`\mu ,\nu ,\mathrm{}`$ are four-dimensional brane indices. The brane metric $`\eta _{\mu \nu }`$ is the four-dimensional flat Minkowski metric with signature $`(,+,+,+)`$. Moreover, $`A(r)=2k|r|`$ where $`k`$ is a positive constant and the fifth dimension $`r`$ runs from $`\mathrm{}`$ to $`\mathrm{}`$. We consider a physical situation such that a single flat 3-brane sits at the origin of the fifth dimension, $`r=0`$, and then ask whether nonabelian gauge fields in a bulk can be localized on the brane only by a gravitational interaction. In this article, we will neglect the back-reaction on the metric from bulk fields, so we do not need to solve five-dimensional Einstein’s equations with the energy-momentum tensor of the bulk fields.

Next we shall review the topological massive nonabelian gauge theories. Such a theory has been first constructed by Yahikozawa and the present author . In those days, our interests have mainly lied in the construction of topological quantum field theories with the topological Higgs mechanism such that there is no metric tensor field in the action and the field equations are merely a flat connection $`F=0`$. Thus in our theory we have not considered the Yang-Mills action $`TrF^2`$, for which we need the metric tensors for contraction of space-time indices. After our work, the topological massive nonabelian gauge theories with usual Yang-Mills kinetic term have been constructed . In the paper at hand, we shall make use of a theory specified to five dimensions of the latter theories.

In five space-time dimensions, the action of the topological massive nonabelian gauge theory is given by

$`S`$ $`=`$ $`{\displaystyle }d^5x[{\displaystyle \frac{1}{4}}\sqrt{g}g^{M_1N_1}g^{M_2N_2}F_{M_1M_2}^aF_{N_1N_2}^a{\displaystyle \frac{1}{48}}\sqrt{g}g^{M_1N_1}g^{M_2N_2}g^{M_3N_3}g^{M_4N_4}`$ (2)
$`\times `$ $`_{M_1M_2M_3M_4}^a_{N_1N_2N_3N_4}^a+{\displaystyle \frac{m}{12}}\epsilon ^{M_1M_2M_3M_4M_5}C_{M_1M_2M_3}^aF_{M_4M_5}^a].`$
Here we have to explain our conventions and notations. Though it is more convenient to use differential forms in order to express various formulas, for later convenience we will make use of the conventional coordinate-dependent notations. For definiteness, we shall take the gauge group to be $`SU(N)`$, with generators $`T^a`$ satisfying

$`[T^a,T^b]`$ $`=`$ $`if^{abc}T^c,`$
$`Tr(T^aT^b)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\delta ^{ab}.`$ (3)
The field strengths $`F`$ and $`H`$ of a gauge field $`A`$ and a 3-form $`C`$ are respectively defined as

$`F_{MN}^a`$ $`=`$ $`_MA_N^a_NA_M^a+gf^{abc}A_M^bA_N^c,`$
$`H_{MNPQ}^a`$ $`=`$ $`4(D_{[M}C_{NPQ]})^a`$ (4)
$`=`$ $`(D_MC_{NPQ})^a(D_NC_{MPQ})^a+(D_PC_{MNQ})^a(D_QC_{MNP})^a,`$
where the covariant derivative $`D`$ is defined in a usual way as $`(D_MC_{NPQ})^a=_MC_{NPQ}^a+gf^{abc}A_M^bC_{NPQ}^c`$. A newly introduced field strength $``$ is defined as $`_{MNPQ}^a=(D_MC_{NPQ})^a+gf^{abc}F_{MN}^bV_{PQ}^c+`$ (cyclic terms) with a 2-form auxiliary field $`V`$. (Henceforth, we set the coupling constant $`g`$ to be 1 except when the coupling constant is needed.) This modified field strength $``$ has been introduced to compensate for the non-invariance of the kinetic term, $`TrH^2`$, under the tensor gauge transformations associated with $`C_{MNP}^a`$ . Note that the $`V_{MN}^a`$ is a non-dynamical field since there is no kinetic term for it in the action. The action (2) is invariant under the gauge transformations

$`\delta A_M^a`$ $`=`$ $`(D_M\theta )^a=_M\theta ^a+f^{abc}A_M^b\theta ^c,`$
$`\delta C_{MNP}^a`$ $`=`$ $`3(D_{[M}\omega _{NP]})^a+f^{abc}C_{MNP}^b\theta ^c,`$
$`\delta V_{MN}^a`$ $`=`$ $`\omega _{MN}^a+f^{abc}V_{MN}^b\theta ^c.`$ (5)
The equations of motion following from the action (2) read

$`[F_{M_1M_2},^{M_1M_2M_3M_4}]^a=0,`$ (6)
$`\left(D_{M_2}(\sqrt{g}F^{M_1M_2})\right)^a+{\displaystyle \frac{1}{6}}\sqrt{g}[C_{M_2M_3M_4},^{M_1M_2M_3M_4}]^a`$
$`+`$ $`{\displaystyle \frac{1}{3}}\left(D_{M_2}(\sqrt{g}[V_{M_3M_4},^{M_1M_2M_3M_4}])\right)^a{\displaystyle \frac{m}{6}}\epsilon ^{M_1M_2M_3M_4M_5}(D_{M_2}C_{M_3M_4M_5})^a=0,`$
$`\left(D_{M_4}(\sqrt{g}^{M_1M_2M_3M_4})\right)^a{\displaystyle \frac{m}{2}}\epsilon ^{M_1M_2M_3M_4M_5}F_{M_4M_5}^a=0.`$
It is easy to show that the topological Higgs mechanism in the nonabelian case occurs in a set of equations of motion (6).

Now using this topological massive nonabelian theory we wish to show that nonabelian gauge field in a five-dimensional bulk is confined to a flat brane. For this, we start by taking the gauge conditions of the symmetries (5)

$`A_r^a(x^M)`$ $`=`$ $`0,`$
$`V_{MN}^a(x^M)`$ $`=`$ $`0.`$ (7)
Furthermore, for simplicity we make an ansatz

$`C_{rMN}^a(x^M)=0.`$ (8)
With these gauge conditions (7) and the ansatz (8), we make the following simple Kaluza-Klein reduction ansatzs for zero-mode

$`A_\mu ^a(x^M)`$ $`=`$ $`a_\mu ^a(x^\lambda )u(r),`$
$`C_{\mu \nu \rho }^a(x^M)`$ $`=`$ $`c_{\mu \nu \rho }^a(x^\lambda )u(r).`$ (9)
where we assume the free equations of motion in four-dimensional flat space-time:

$`^\mu \overline{f}_{\mu \nu }^a=^\mu \overline{h}_{\mu \nu \rho \sigma }^a=0,`$ (10)
with the definitions being $`\overline{f}_{\mu \nu }^a=2_{[\mu }a_{\nu ]}^a`$ and $`\overline{h}_{\mu \nu \rho \sigma }^a=4_{[\mu }c_{\nu \rho \sigma ]}^a`$.

Even with these simple ansatzs, it is difficult to find the zero-mode solution $`u(r)`$ owing to the non-linearity of the equations of motion. However, as in the analysis of the graviton where the linear fluctuations $`h_{\mu \nu }`$ around the flat metric $`\eta _{\mu \nu }`$ are studied at the lowest level of approximation, it is sufficient to take account of the linear equations of the fields, which implies that we look for a solution at the lowest level of the coupling constant $`g`$. The linear equations of motion for the fields reduce to

$`_{M_2}\left(\sqrt{g}g^{M_1N_1}g^{M_2N_2}\overline{F}_{N_1N_2}^a\right){\displaystyle \frac{m}{6}}\epsilon ^{M_1M_2M_3M_4M_5}_{M_2}C_{M_3M_4M_5}^a`$ $`=`$ $`0,`$
$`_{M_4}(\sqrt{g}g^{M_1N_1}g^{M_2N_2}g^{M_3N_3}g^{M_4N_4}\overline{H}_{N_1N_2N_3N_4})m\epsilon ^{M_1M_2M_3M_4M_5}_{M_4}A_{M_5}^a`$ $`=`$ $`0,`$ (11)
where $`\overline{F}_{MN}^a=2_{[M}A_{N]}^a`$ and $`\overline{H}_{MNPQ}^a=4_{[M}C_{NPQ]}^a`$. Then, using Eqs. (1) and (9), Eq. (11) takes the forms

$`a^{a\mu }_r\left(\text{e}^{A(r)}_ru(r)\right){\displaystyle \frac{m}{6}}\epsilon ^{\mu \nu \rho \sigma }c_{\nu \rho \sigma }^a_ru(r)=0,`$ (12)
$`\text{e}^{A(r)}^\mu a_\mu ^a_ru(r){\displaystyle \frac{m}{6}}\epsilon ^{\mu \nu \rho \sigma }_\mu c_{\nu \rho \sigma }^au(r)=0,`$ (13)
$`c^{a\mu \nu \rho }_r\left(\text{e}^{A(r)}_ru(r)\right)m\epsilon ^{\mu \nu \rho \sigma }a_\sigma ^a_ru(r)=0,`$ (14)
$`\text{e}^{A(r)}_\rho c^{a\mu \nu \rho }_ru(r)m\epsilon ^{\mu \nu \rho \sigma }_\rho a_\sigma ^au(r)=0.`$ (15)
Note that up to gauge group index $`a`$, these equations have the same forms as those in the case of abelian gauge field . As the abelian case, since we can regard Eqs. (13), (15) as the gauge conditions in four-dimensional space-time, the equations which we have to solve are Eqs. (12), (14). From the two equations, we can derive a differential equation to $`u(r)`$:

$`\left(_ru(r)\right)^2m^2u^2(r)=0.`$ (16)
Provided that after $`Z_2`$ orbifolding with respect to the fifth dimension we impose an even reflection symmetry on $`u(r)`$, a general solution to (16) is given by

$`u(r)=c\text{e}^{\pm m\epsilon (r)r},`$ (17)
where $`c`$ is an integration constant and $`\epsilon (r)`$ is the step function defined as $`\epsilon (r)=\frac{r}{|r|}`$ and $`\epsilon (0)=0`$. Moreover, imposing the boundary conditions such that $`u(\pm \mathrm{})=0,u(0)=u_0`$, a special solution takes the form

$`u(r)=u_0\text{e}^{m\epsilon (r)r}+𝒪(g).`$ (18)
Here in order to emphasize that this solution is a solution up to the leading order of the gauge coupling constant $`g`$, we have put the term $`𝒪(g)`$.

We are now ready to check that this zero-mode solution of the gauge field as well as the 3-form potential leads to a normalizable mode and the localization on a brane. To do so, plugging Eqs. (7), (8) and (9) into the starting action (2), we have

$`S^{(0)}`$ $`=`$ $`{\displaystyle }d^4x{\displaystyle _{\mathrm{}}^{\mathrm{}}}dr[{\displaystyle \frac{1}{4}}F_{\mu \nu }^aF^{a\mu \nu }{\displaystyle \frac{1}{2}}a_\mu ^aa^{a\mu }\text{e}^A(_ru)^2`$ (19)
$``$ $`{\displaystyle \frac{1}{48}}\text{e}^{2A}H_{\mu \nu \rho \sigma }^aH^{a\mu \nu \rho \sigma }{\displaystyle \frac{1}{12}}c_{\mu \nu \rho }^ac^{a\mu \nu \rho }\text{e}^A(_ru)^2],`$
where the topological term has dropped from the above action after performing integration by parts over $`r`$. First, let us focus our attention to the first plus second terms in (19). The calculation of the $`r`$-integrals leads to

$`S_1`$ $``$ $`{\displaystyle d^4x_{\mathrm{}}^{\mathrm{}}𝑑r\left[\frac{1}{4}F_{\mu \nu }^aF^{a\mu \nu }\frac{1}{2}a_\mu ^aa^{a\mu }\text{e}^A(_ru)^2\right]}`$ (20)
$`=`$ $`{\displaystyle }d^4x[{\displaystyle \frac{1}{4}}\{{\displaystyle \frac{u_0^2}{m}}(_\mu a_\nu ^a_\nu a_\mu ^a)^2+{\displaystyle \frac{4u_0^3}{3m}}gf^{abc}(_\mu a_\nu ^a_\nu a_\mu ^a)a^{b\mu }a^{c\nu }`$
$`+`$ $`{\displaystyle \frac{u_0^4}{2m}}g^2f^{abc}f^{ade}a_\mu ^ba_\nu ^ca^{d\mu }a^{e\nu }\}{\displaystyle \frac{1}{2}}{\displaystyle \frac{m^2u_0^2}{m+k}}a_\mu ^aa^{a\mu }],`$
where the coupling constant $`g`$ was recovered. In order to transform the free kinetic terms to a canonical form, let us redefine the field and the coupling constant as

$`{\displaystyle \frac{u_0}{\sqrt{m}}}a_\mu ^a`$ $``$ $`a_\mu ^a,`$
$`{\displaystyle \frac{2\sqrt{m}}{3}}g`$ $``$ $`g.`$ (21)
Consequently, the action (20) reduces to

$`S_1`$ $`=`$ $`{\displaystyle }d^4x[{\displaystyle \frac{1}{4}}\{(_\mu a_\nu ^a_\nu a_\mu ^a)^2+2gf^{abc}(_\mu a_\nu ^a_\nu a_\mu ^a)a^{b\mu }a^{c\nu }`$ (22)
$`+`$ $`{\displaystyle \frac{9}{8}}g^2f^{abc}f^{ade}a_\mu ^ba_\nu ^ca^{d\mu }a^{e\nu }\}{\displaystyle \frac{1}{2}}{\displaystyle \frac{m^3}{m+k}}a_\mu ^aa^{a\mu }].`$
Let us note that the coefficient in front of the third term of order $`𝒪(g^2)`$ is not 1 but $`\frac{9}{8}`$, which makes impossible to transform the parts except a mass term to the square form. This is because the solution (18) is a solution holding only up to $`𝒪(g)`$, thereby meaning that terms of $`𝒪(g^2)`$ in the action are ambiguous. To determine the correct coefficient in front of the third term, we require the gauge invariance for the action except a mass term. As a result, the action (22) reads

$`S_1={\displaystyle d^4x\left[\frac{1}{4}f_{\mu \nu }^af^{a\mu \nu }\frac{1}{2}\frac{m^3}{m+k}a_\mu ^aa^{a\mu }\right]},`$ (23)
where $`f_{\mu \nu }^a`$ is the four-dimensional field strength of the gauge field $`a_\mu ^a`$, $`f_{\mu \nu }^a=_\mu a_\nu ^a_\nu a_\mu ^a+gf^{abc}a_\mu ^ba_\nu ^c`$.

Next let turn to the third term in (19), that is, $`\frac{1}{48}d^5x\text{e}^{2A}H_{\mu \nu \rho \sigma }^aH^{a\mu \nu \rho \sigma }`$. Since we have the expression $`H_{\mu \nu \rho \sigma }^a=_\mu c_{\nu \rho \sigma }^au(r)+gf^{abc}a_\mu ^bc_{\nu \rho \sigma }^cu^2(r)+`$ (cyclic terms), the $`r`$-integral in front of the free kinetic term for $`c_{\mu \nu \rho }^a`$ becomes

$`I_3{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑r\text{e}^{2A}u^2(r)={\displaystyle \frac{u_0^2}{m2k}}.`$ (24)
when $`m2k>0`$, whereas it diverges when $`m2k0`$. As in abelian gauge field, we assume the relation $`m2k0`$, by which a 3-form potential $`c_{\mu \nu \rho }^a`$ is not localized on a brane, but lives in a bulk away from the brane . Accordingly, the brane action is given by the action (23). In addition to the relation $`m2k0`$, we have to require the massless condition of the gauge field, which becomes $`\frac{m^3}{m+k}1`$. These conditions are simply satisfied by taking $`km`$ as in the case of abelian gauge field, that is, ’$`\mathrm{𝑝ℎ𝑜𝑡𝑜𝑛}`$’. Hence we have shown that like abelian gauge field, nonabelian gauge field is also localized on a brane by a gravitational interaction.

At this stage, we should comment on one subtlety. Even if we have shown that gauge fields with almost vanishing mass can be localized on a brane by a gravitational interaction, exactly massless gauge fields cannot be localized on the brane by the mechanism mentioned so far. This fact can be understood by looking at the normalized zero-mode in a flat space-time, $`\widehat{u}(r)=\sqrt{m}\text{e}^{m\epsilon (r)r}`$, which becomes vanishing when $`m`$ is exactly zero. This situation is the same as that of fermions where only fermions with mass of a ’kink’ profile can be localized on the brane . This problem does not give us any problem to massless gauge fields since in the case of the abelian gauge theory the zero-mass case is simply the limiting case of the finite mass theory . However, it is well-known that in the nonabelian gauge theory and the gravitational theory, there is a discrete difference between the theory with zero-mass and the theory with finite mass, no matter how small as compared to all external momenta. In particular, this problem is serious in the gravitational theory since experiments of the bending of light rays near the sun and the perihelion movement of Mercury are distinctly different for the zero-mass graviton and the graviton with infinitesimally small mass. Of course, experiments force us to select the zero-mass theory.

This situation is more subtle for the gauge theory compared to the gravitational theory by the following two reasons. In the case of the gravity, the discontinuity appears at the tree level, while in the case of the nonabelian gauge theory it occurs at the loop level, so the corrections might be very small in this case. The second reason is that although we theoretically regard the gluons as exact massless particles, the gluons with a mass as large as a few MeV may not be precluded experimentally . Thus if the gluons in fact have such a small mass, we can apply the present localization mechanism to nonabelian gauge fields without conflicting experiments.

In conclusion, we have applied a new localization mechanism developed for abelian gauge fields to the case of nonabelian gauge fields. Our mechanism fully utilizes the topologically massive nonabelian gauge theories and is very similar to that of fermions in the sense that in the both mechanisms the zero-modes share the same form and the presence of the mass term of a ’kink’ profile plays an essential role for trapping the zero-modes of the bulk fields on a flat Minkowski brane. It is worthwhile to stress that abelian gauge fields and nonabelian gauge fields in a bulk are localized on a brane by the same mechanism. We have also pointed out one subtlety associated with massless nonabelian gauge fields.


text:
# The order-disorder character of the (3×3) to (√3×√3)R30∘ phase transition of Sn on Ge(111)

## Abstract

The $`\alpha `$-phase of Sn/Ge(111) has been investigated from 120 K up to 500 K, using synchrotron radiation core level photoemission. By means of photoelectron diffraction experiments, we verified that the rippled structure of the low temperature $`(3\times 3)`$-phase is preserved in the $`(\sqrt{3}\times \sqrt{3})`$R30-phase at room temperature, thus confirming the order-disorder character of the phase transition. We also found that at least two components are present in the Sn 4d core level spectra up to 500 K, i.e about 300 K above the onset of the transition from the low temperature $`(3\times 3)`$-phase to the $`(\sqrt{3}\times \sqrt{3})`$R30-phase, thus excluding the occurence of any displacive transition.

The system formed by one third of a monolayer (ML) of Sn adatoms atop the T4 sites of the Ge(111) surface ($`\alpha `$-phase) exhibits a transition from a low temperature (LT) $`(3\times 3)`$ phase to a high temperature (HT) $`(\sqrt{3}\times \sqrt{3})`$R30 phase. This transition is gradual, reversible and has a critical temperature T$`{}_{c}{}^{}`$ 220 K. The phase transition was first attributed to the manifestation of a commensurate surface charge density wave state in the LT phase stabilized by many body effects, in particular by electron correlation. Later on, an order-disorder model, invoking dynamic fluctuations as microscopic driving mechanism, has been proposed to explain the published data. The research then focused on the role of surface defects and the determination of the surface structure in an attempt to discriminate between the two interpretations.

At present, the $`(3\times 3)`$ structure has been determined with a good agreement between surface X-ray diffraction (SXRD) and photoelectron diffraction (PED) measurements. At low temperature, the Sn adatoms occupy inequivalent $`T_4`$ sites, with one Sn adatom (out of three within the unit cell) protruding above the surface. This inequivalency is also confirmed by photoemission experiments, where at least two components (with intensity ratio close to 1:2) are required to fit the Sn 4d spectrum properly.

For the structure of the $`(\sqrt{3}\times \sqrt{3})`$R30 phase at room temperature (RT), only SXRD measurements are available. They slightly favor a structure with equivalent Sn adatoms, thus suggesting that the transition is indeed due to a pseudo-Jahn-Teller effect, i.e. it would be a displacive phase transition. This result is hardly conciliating with a very recent He diffraction study of the $`(3\times 3)`$ order parameter, which indicates the occurence of a 3-state Potts order-disorder phase transition at $`T_c220`$ K. Moreover, the hypothesis of a displacive transition at 220 K is also not compatible with the core level photoemission experiments, where the inequivalency between the Sn atoms seems to remain also in the $`(\sqrt{3}\times \sqrt{3})`$R30 phase, since the Sn 4d lineshape does not change significantly from $``$ 100-150 K up to room temperature. This discrepancy urges for PED structural measurements of the RT phase.

The hypothesis of a displacive transition, being recently suggested to be driven by a phonon softening, asks for experiments in a more extended temperature range. To our knowledge no data have been published at temperatures higher than 300 K, while it is well known that the displacive phase transitions are always accompanied by an intermediate disordered stage around the critical temperature, manifesting the change in surface corrugation only at a temperature substantially higher than Tc. It is therefore of utmost importance to have a clear indication of the order-disorder character of the transition around Tc and to know if the two components observed in the Sn 4d core level spectra at LT and RT are also present far from Tc, up to the highest temperature where this phase is stable ($``$ 550 K).

By means of photoelectron diffraction (PED) measurements, we have already shown that in the $`(3\times 3)`$ phase the observed splitting of the Sn 4d core levels in two components derives from adatoms with different bonding configuration. A vertical ripple of 0.3 Å between the Sn adatoms has been estimated by simulations, with one adatom out of three that protrudes above the surface. In particular, the adatom at the higher height level is associated to the Sn 4d component with the larger binding energy. The three nearest-neighbor Ge atoms follow the Sn adatom vertical distortion, i.e. their bond angle and length remain almost unchanged. This rippled structure is in good agreement with both the SXRD experiments and recent density- functional theory (DFT) calculations for the LT $`(3\times 3)`$-phase.

In this paper, we report on the evolution of the Sn 4d core levels and on their PED patterns at different temperatures. The measurements were performed at the ELETTRA Synchrotron facility in Trieste, in the ultra high vacuum end-stations of the ALOISA and SuperESCA beamlines. The Sn 4d core levels photoemission spectra presented here were measured at the SuperESCA beamline with an overall energy resolution better than 120 meV. The PED experiments were performed at the ALOISA beamline.

Details about the sample preparation and PED experiment at ALOISA are given elsewhere. At the SuperESCA beamline, the $`\alpha `$-phase was prepared by dosing 1/3 ML of Sn on a clean Ge(111)-c$`(2\times 8)`$ surface kept at RT. The optimum coverage and annealing treatment for the $`(\sqrt{3}\times \sqrt{3})`$R30 and $`(3\times 3)`$ phases was checked by optimizing the intensity of the LEED $`(3\times 3)`$ extra spots of the LT phase and the intensity of the Sn induced surface band at the Fermi level in the photoemission spectra. This feature is indeed very sensitive to the quality of the $`(\sqrt{3}\times \sqrt{3})`$R30 phase. We have investigated also the $`(2\times 2)`$ phase of Sn/Ge(111), obtained by dosing $``$ 0.2 ML of Sn and annealing at 450 K. The temperature stability during measurements was better than $`\pm `$ 10 K.

In Fig. 1, we show the Sn 4d core level photoemission spectra taken at room temperature in normal emission geometry for photon energies between 120 and 300 eV. In these experimental conditions, the photon energy dependence of the photoemission spectra (after normalization to the photoemission cross-section) mainly reflects the vertical distance between the Sn adatom and the Ge atom directly underneath, thus enhancing the sensitivity to the vertical distortions of the adatom layer. In fact, the Sn vertical ripple is the main structural parameter involved in the phase transition, since it is expected to disappear in the high temperature phase if a displacive transition takes place. A detailed description of the Sn 4d core level analysis in terms of two spin-orbit split doublets is given in Ref. , where high resolution spectra have been studied to determine the fitting parameters used for the successive analysis of the PED data taken in the LT phase. Also in the $`(\sqrt{3}\times \sqrt{3})`$R30 phase at 300 K, the Sn 4d core levels have a shape that cannot be fitted with a single spin-orbit split doublet. At least one additional doublet is needed to improve appreciably the quality of the fit (see the inset panel of Fig. 1).

The photon energy dependence of the intensity ratio $`I_B/I_A`$ between the minority $`B`$ and the majority $`A`$ components of the Sn 4d5/2 core levels at 130 K and 300 K is shown in Fig. 2. This presentation of the PED data allows a direct evaluation of any change of the Sn vertical ripple, in particular, the disappearance of the ripple would yield a constant intensity ratio without any energy modulation. In addition, the comparison between the intensity ratio below and above the transition is not affected by the sistematic errors of extracting the anisotropy $`\chi `$-function. These errors can be large around 130 eV, because of the pronounced Cooper minimum. The error bars in Fig. 2 represent the maximum error in the intensity ratio, as obtained by the fitting procedure of ref. , and it mainly arises from the uncertainty on the energy shift between the $`A`$ and $`B`$ components and their linewidth. The intensity ratio $`I_B/I_A`$ below and above the phase transition is very similar, confirming that, also in the $`(\sqrt{3}\times \sqrt{3})`$R30 phase at RT, there are two inequivalent types of Sn adatoms with different bonding geometry (different vertical distances between Sn and the underneath Ge atom). Thus, the rippled structure is maintained up to RT. This result differs from that obtained by SXRD experiments (where the surface corrugation was claimed to change from the $`(3\times 3)`$ to the $`(\sqrt{3}\times \sqrt{3})`$R30 phase), and indicates that the phase transition at T$`{}_{c}{}^{}`$ 220 K has an order-disorder character.

Since the two spectral components are unambiguously assigned to two adsorption sites with different Sn heights, we can detect any structural modification of the $`(\sqrt{3}\times \sqrt{3})`$R30 phase by simply measuring the Sn 4d spectrum as a function of the temperature. In particular, an eventual mergence of the two spectral components in a single one would be the fingerprint of a displacive transition. We have reported in Fig. 3 the lineshapes of the Sn 4d core level spectra of 1/3 ML of Sn/Ge(111) measured from 120 K to 500 K with a photon energy of 198 eV and an emission angle of 55 from the surface normal. These spectra have been fitted to Doniach Sunjic (DS) doublets convoluted with a Gaussian. The Gaussian takes into account both the instrumental broadening (120 meV) and the temperature dependent phonon broadening, that is much larger. The free parameters were the energy position of the doublets, their intensities, the Gaussian full width at half maximum (FWHM) \[assumed to be the same for all the doublets\], the asymmetry of the DS line, and the branching ratio. The spin- orbit splitting was set to 1.02 eV (according to Refs.  and a polynomial function was used for the background caused by the low energy tail of the Ge 3d core levels at about 29 eV of binding energy.

Besides the $`A`$ and $`B`$ components, a third weak doublet was required to fit properly the $`(3\times 3)`$ phase spectra at 120 K. This component is attributed to the Sn atoms surrounding Ge substitutional impurities in the $`(3\times 3)`$ overlayer, according to the measurements of Uhrberg et al.. We observed that the third component intensity decreases by increasing the temperature up to disappearance at 420 K. This is possibly due to limits of the fitting procedure in detecting such a faint feature when the main components become increasingly broader. The energy shift between the $`A`$ and $`B`$ components does not change between 120 K and 500 K (see the upper panel of Fig. 4), while the corresponding gaussian width gradually broadens from $`0.3`$ eV up to $`0.5`$ eV. The Sn 4d core level lineshape at 500 K is more symmetric and can be fitted to a single DS doublet, thus yielding a Gaussian FWHM of $`0.6`$ eV, a value much larger than that of the components used to fit the low temperature spectra.

In order to discriminate between the model with a single component at 500 K and that with $`A`$ and $`B`$ components, we have followed the temperature dependence of the Gaussian FWHM of these components in the $`(3\times 3)`$, in the $`(\sqrt{3}\times \sqrt{3})`$R30 and in the $`(2\times 2)`$ phases of Sn on Ge(111). The corresponding temperature dependences are reported in Fig. 4. The values obtained for the $`A`$ and $`B`$ components of the $`\alpha `$-phase and for the $`(2\times 2)`$ phase show a linear increase with the temperature, with a slope of $`(0.5\pm 0.02)10^3`$ eV K-1 and $`(0.32\pm 0.03)10^3`$ eV K-1, respectively. This result is consistent with the linear increase ($`0.210^3`$ eV K-1) observed for the surface components of the 3d core levels of the clean Ge(111) surface in a wider temperature range. The fit of the $`\alpha `$-phase core levels with a single component yields a gaussian width which does not change with the temperature, while the quality of the fit declines as the temperature is lowered (being definitely insufficient already at 300 K). If the 500 K spectra contained a single component, there would be a jump of $`0.2`$ eV of the FWHM between 300 K and 500 K, with about a threefold increase in the temperature derivative of the core level width. This implies that the thermal broadening of the Sn 4d core levels would be strongly non linear, in contrast to the data of the Sn $`(2\times 2)`$ phase, the clean Ge(111) surface and other semiconductors surfaces. As a consequence, the temperature dependence of the FWHMs reported in Fig. 4 strongly supports a deconvolution of the Sn 4d core level spectra with at least two DS doublets up to 500 K.

In conclusion, we have measured the temperature and energy dependence of the Sn 4d core levels of Sn/Ge(111). By means of energy dependent PED, we found that the vertical distortion, obtained for the two kinds of Sn adatoms in the LT $`(3\times 3)`$ phase, is also present in the $`(\sqrt{3}\times \sqrt{3})`$R30 phase at RT, i.e. the transition at $``$ 220 K has an order-disorder character. In addition, we have shown that the two components of the photoemission spectra, reflecting the inequivalence of the Sn adsorption sites, are present up to 500 K. This result leads us to discard the occurrence of a displacive character of the transition even at temperatures much higher than T$`{}_{c}{}^{}=220`$ K.

This work was partly supported by MURST cofin99 (Prot. 9902332155 and Prot. 9902112831), by Regione Friuli-Venezia Giulia 98 and by INFM PAIS-F99.