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# Hyperbolic localization of intersection cohomology

## 1. Introduction

Let $`X`$ be a normal variety over an algebraically closed field $`k`$ on which the the multiplicative group $`T=G_m(k)=k^\times `$ acts algebraically. Let $`F=X^T`$ be the variety of fixed points, with connected components $`F_1,\mathrm{},F_r`$. Define varieties

$$X_k^+=\{xX\underset{t0}{lim}txF_k\},\text{ and}$$

$$X_k^{}=\{xX\underset{t\mathrm{}}{lim}txF_k\}.$$

Let $`X^+`$ (respectively $`X^{}`$) be the disjoint (disconnected) union of all the $`X_k^+`$ (resp. $`X_k^{}`$), and let $`f^\pm :FX^\pm `$ and $`g^\pm :X^\pm X`$ restrict to the inclusion maps on each component $`F_k`$, $`X_k^\pm `$.

If $`k=`$, let $`𝒟(X)`$ the category of constructible complexes of sheaves of on $`X`$ over a commutative ring $`R`$ of finite global dimension. Otherwise, let $`𝒟(X)`$ be either the derived category of étale sheaves on $`X`$ with torsion coefficients prime to the characteristic of $`k`$, or one of the categories $`D^b(X;_l)`$, $`D^b(X;_l)`$, $`D^b(X;\overline{_l})`$ obtained from them by standard limiting arguments.

Define hyperbolic localization functors $`()^!,()^!:𝒟(X)𝒟(F)`$ by

$$S^!:=(f^+)^!(g^+)^{}S,$$

$$S^!:=(f^{})^{}(g^{})^!S.$$

In general these functors are different; our first theorem gives sufficient conditions for them to be canonically isomorphic. Let $`\mu :T\times XX`$ be the map defining the action. Call an object $`S𝒟(X)`$ “weakly equivariant” if $`\mu ^{}SLS`$ for some locally constant sheaf $`L`$ on $`T`$. Equivariant perverse sheaves, and more generally objects of the equivariant derived category of \[BL\] give rise to weakly equivariant objects. Call $`S`$$`T`$-constructible” if it is constructible with respect to a stratification by $`T`$-invariant subvarieties. Weakly equivariant objects are $`T`$-constructible, but not conversely.

###### Theorem 1.

There is a natural morphism $`\iota _S:S^!S^!`$, which is an isomorphism if $`S`$ is weakly equivariant, or if $`k=`$, $`R`$ is an algebraically closed field, and $`S`$ is $`T`$-constructible.

Roughly speaking, the stalks of $`S^!`$ and $`S^!`$ can be thought of as the local cohomology of $`S`$ at a point of $`F`$ with mixed supports — closed supports in the directions flowing into $`F`$, and compact supports in directions flowing away. These two kinds of supports come from limiting operations that do not in general commute; the content of Theorem 1 is that they do commute for weakly equivariant objects.

To motivate our second result, consider the case where $`X`$ is a smooth projective variety over $``$. Then the $`F_k`$ are all smooth, and the varieties $`X_k^+`$, $`X_k^{}`$ are affine bundles over them, by \[B-B\]. Localizing the constant sheaf $`_X`$ gives

$$(_X)^!=\underset{k=1}{\overset{r}{}}_{F_k}[n_k],$$

where $`n_k=dim_{}X_k^+dim_{}F_k`$, and there is an isomorphism between $`H^{}(X;)`$ and the hypercohomology

$$^{}((_X)^!)=\underset{k=1}{\overset{r}{}}H^{n_k}(F_k;).$$

Kirwan \[Ki\] generalized this result to the intersection cohomology $`IH^{}(X;)`$ of a singular complex projective variety $`X`$: there is a non-canonical isomorphism $`IH^{}(X;)^{}(F;\mathrm{𝐈𝐂}^{}(X;)^!)`$. Here $`\mathrm{𝐈𝐂}^{}(X)`$ is the intersection cohomology complex of \[GM\].

To make Kirwan’s result useful, one wants to be able to compute $`\mathrm{𝐈𝐂}^{}(X)^!`$. Our second result shows that it is as simple as one could hope. We assume either that $`k=`$ or else that $`k=\overline{𝔽_q}`$ and $`X`$ is obtained by extension of scalars from a variety defined over $`𝔽_q`$. If $`k=`$, we take $`\mathrm{𝐈𝐂}^{}(X)=\mathrm{𝐈𝐂}^{}(X;)`$; if $`k=\overline{𝔽_q}`$, we take $`\mathrm{𝐈𝐂}^{}(X)`$ to be the intersection cohomology complex in $`D^b(X;\overline{_l})`$.

###### Theorem 2.

There is an isomorphism

$$\mathrm{𝐈𝐂}^{}(X)^!\underset{j=1}{\overset{n}{}}\mathrm{𝐈𝐂}^{}(Y_j,_j)[d_j]$$

where $`Y_1,\mathrm{},Y_n`$ are irreducible closed subvarieties of $`F`$, each $`_j`$ is a local system on a Zariski open subset of the smooth locus of $`Y_j`$, and $`d_j`$.

###### Remarks.

(1) To see why either weak equivariance or constructibility is needed in Theorem 1, take $`X=^2`$ with the action $`t(x,y)=(tx,t^1y)`$, and let $`S=i_{}_{}`$, where $`i:X`$ is the diagonal embedding. Then $`S^!=j^!S`$ and $`S^!=j^{}S`$, where $`j`$ is the inclusion of the origin into $`X`$.

(2) Replacing the torus action by the opposite action interchanges $`X^+`$ and $`X^{}`$, so the hyperbolic localization functors are sent to their Verdier duals. Thus up to duality there is only one hyperbolic localization functor for weakly equivariant sheaves.

(3) In \[GM\] Goresky and MacPherson calculated local contributions to the Lefschetz fixed-point formula for a self-map (or more generally a correspondence) acting on a complex of sheaves on a singular space. They assume that the map is “weakly hyperbolic”, which ensures that attracting and repelling directions can be defined. They give a sufficient condition (Proposition 9.2) for $`^{}(S^!)`$ and $`^{}(S^!)`$ to be isomorphic. This condition is satisfied when the self-map comes from a torus action; this can be used to show that Theorem 1 holds when $`k=`$ and $`S`$ is $`T`$-constructible, without the assumption on the coefficient ring $`R`$.

(4) Theorem 2 is proved by using Theorem 1 to show that hyperbolic localization preserves purity of weakly equivariant mixed Hodge modules or mixed $`l`$-adic sheaves. Similar arguments were used in \[Gi\] \[BrM\] in the special case when the action is completely attracting near a fixed point component; in this case Theorem 1 is obvious.

(5) Except for this special case, it is not possible to replace hyperbolic localization by ordinary restriction in Theorem 2. For instance, take $`X`$ to be a quadric cone in $`M=^4`$ with equation $`xyzw=0`$. Let the $`T`$-action be given by $`t(x,y,z,w)=(tx,t^1y,z,w)`$. Then $`FX=X^T`$ is the union of two lines, and $`\mathrm{𝐈𝐂}^{}(X)|_F_{FX}_0[2]`$, which is not a sum of intersection cohomology sheaves. A further computation with this example also shows that none of the functors $`(f^+)^!`$, $`(f^{})^{}`$, $`(g^+)^{}`$, $`(g^{})^!`$ preserve purity individually.

Finally, we mention a useful property of hyperbolic localization, which is used together with Theorem 2 in \[BiBr\] to give inequalities on Kazhdan-Lusztig polynomials. We restrict to the case of sheaves on $`X`$ with coefficients in a field.

###### Proposition 3.

Hyperbolic localization preserves local Euler characteristics: given any point $`xF`$ and any $`S𝒟(X)`$, we have

$$\chi (S_x)=\chi ((S^!)_x).$$

Since $`(f^+)^!=𝔻_F(f^+)^{}𝔻_{X^+}`$, this follows from the fact that Verdier duality preserves local Euler characteristics. One way to see this is to note that $`𝒟(X)`$ is generated as a triangulated category by (shifted) intersection cohomology sheaves, and then use the fact that $`𝔻_X\mathrm{𝐈𝐂}^{}(Y)=\mathrm{𝐈𝐂}^{}(Y)`$. For varieties over $``$, this can also be deduced from \[Su\] or \[KS, Exercise IX.12\].

### Acknowledgments

The author would like to thank Mikhail Finkelberg, David Massey, Ivan Mirković, Kari Vilonen, and the referees for useful comments and corrections, and Nick Katz for pointing out that Lemma 7 does not hold in the étale setting.

## 2. Reduction to the affine case

The following is an easy consequence of proper base change.

###### Lemma 4.

The hyperbolic localization functors commute with pullbacks by inclusions of open $`T`$-invariant subvarieties and with pushforwards by inclusions of closed $`T`$-equivariant subvarieties.

To construct the natural morphism $`\iota _S:S^!S^!`$ of Theorem 1, it is enough to construct its restriction to each component of $`F`$. By removing all the components of $`F`$ but one and using the lemma, we can assume that $`X`$ has only one fixed point component. In this case $`X^+X^{}=F`$, so the maps $`f^\pm `$, $`g^\pm `$ form a Cartesian square. Let $`S^+=(g^+)_{}(g^+)^{}S`$, and let $`\beta :SS^+`$ be the adjunction morphism. Base change gives a natural isomorphism

$`(S^+)^!`$ $`=`$ $`(f^{})^{}(g^{})^!(g^+)_{}(g^+)^{}S`$
$``$ $`(f^{})^{}(f^{})_{}(f^+)^!(g^+)^{}SS^!`$
Define $`\iota _S`$ to be the composition of (2) with $`\beta ^!:S^!(S^+)^!`$. It is compatible with the pullback by the inclusion $`j:UX`$ of a Zariski open subset: we have $`j^{}\iota _S\iota _Uj^{}`$.

###### Lemma 5.

Any normal $`T`$-variety defined over an algebraically closed field has a covering by $`T`$-invariant affine open subvarieties, each of which is equivariantly isomorphic to a $`T`$-invariant closed subvariety of an affine space $`𝔸^N`$ on which $`T`$ acts linearly.

This is a consequence of results of Sumihiro \[Sum\]. (We remark that this is the only place where we need that $`k`$ is algebraically closed. Theorem 1 will still hold, for instance, when $`X`$ is a real variety for which the conclusion of Lemma 5 can be verified directly – for instance, if $`X`$ is a $`T_{max}()`$-invariant subvariety of a real flag variety $`G()/B()`$, where $`G`$ is a semisimple algebraic group, $`B`$ is a Borel subgroup, and $`T_{max}`$ is a maximal split torus, all defined over $``$).

As a result, it is enough to prove Theorem 1 when $`X=𝔸^N`$ with a linear $`T`$-action.

## 3. Comparing pushforwards and pullbacks

The proofs of Theorem 1 and Theorem 2 use the following technical lemma. Suppose that $`T=G_m`$ acts trivially on a variety $`Y`$, and that $`WY`$ is a $`T`$-equivariant line bundle, provided with an equivariant splitting $`W=W_+W_{}`$ so that all the weights of $`W_+`$ are larger than all the weights of $`W_{}`$.

Let $`E=(W)(W_+)`$ and $`B=(W_{})`$, where $`(W)`$ denotes the projective bundle of lines in $`W`$. Then $`T`$ acts on $`E`$ so that $`lim_{t0}txB`$ for all $`xE`$. The resulting map $`p:EB`$ is given by $`p([w_+,w_{}])=[0,w_{}]`$. Let $`\varphi :BY`$ and $`i:BE`$ be the projection and inclusion maps, respectively.

###### Lemma 6.

There are natural isomorphisms in $`𝒟(Y)`$

$$\varphi _{}p_{}S\varphi _{}i^{}S,\varphi _{}p_!S\varphi _{}i^!S$$

for any weakly equivariant object $`S𝒟(E)`$.

The proof is in §5.

If $`W_{}`$ is one-dimensional, we get the following special case: $`\varphi :BY`$ is an isomorphism, and $`E`$ is an equivariant vector bundle over $`B`$ with only positive weights. This case of the lemma is well-known; see \[Sp\], for instance. The general case when the action on $`B`$ is non-trivial seems to be new, and is essential for the proof of Theorem 1.

This gives another description of hyperbolic localization. With the notation from the introduction, define projection maps $`\pi ^\pm :X^\pm F`$ by $`\pi ^+(x)=lim_{t0}tx`$ and $`\pi ^{}(x)=lim_t\mathrm{}tx`$.

Adjunction gives natural morphisms

$$(\pi ^{})_{}(f^{})^{},\text{and}(f^+)^!(\pi ^+)_!.$$

They are isomorphisms for weakly equivariant objects, by Lemma 5 and special case of Lemma 6. Thus for any weakly equivariant $`S𝒟(X)`$ there are natural isomorphisms

(1)
$$S^!(\pi ^+)_!(g^+)^{}S,S^!(\pi ^{})_{}(g^{})^!S.$$

## 4. Proof of Theorem 1

First note that if Theorem 1 holds for two objects $`S,S^{}𝒟(X)`$, then it holds for the cone of any morphism $`SS^{}`$. The following lemma then allows us to reduce the proof of Theorem 1 to the case that $`S`$ is weakly equivariant.

###### Lemma 7.

If $`k=`$ and $`R`$ is an algebraically closed field, then the smallest full subcategory of $`𝒟(X)=D^b(X;R)`$ containing all weakly equivariant objects which is closed under taking cones is the category of $`T`$-constructible objects.

###### Proof.

One can show that any local system on a $`T`$-invariant variety has a filtration with weakly equivariant subquotients, by using the Jordan decomposition of the monodromy generated by the action of a loop in $`T`$. The lemma follows from this by induction on strata. ∎∎

Thus we can assume that $`S`$ is weakly equivariant. By the discussion of Section 2 we can also assume that $`X=V`$ is a $`k`$-vector space with a linear $`T`$-action. If $`V=V_+V_{}V_0`$ is the decomposition of $`V`$ into subspaces where $`T`$ acts with positive, negative, and zero weights, respectively, then we have $`F=V_0`$, $`X^+=V_+V_0`$, and $`X^{}=V_{}V_0`$.

Let $`j`$ be the inclusion of the complement of $`X^+`$ into $`X`$, and consider the triangle containing the morphism $`\beta `$ from §2:

$$j_!j^{}SS\stackrel{\beta }{}(g^+)_{}(g^+)^{}S\stackrel{[1]}{}.$$

To show that $`\iota _S`$ is an isomorphism it will be enough to show that $`(j_!j^{}S)^!`$ vanishes, since then $`\beta ^!`$ is an isomorphism. Thus we can assume that $`S|_{X^+}=0`$.

Let $`\overline{X}=V_0\times ((V_+V_{}k)(V_{}))`$; define a $`T`$-action by letting $`T`$ act trivially on the summand $`k`$. There is an equivariant open embedding $`b:X\overline{X}`$ given by $`b(v_+,v_{},v_0)=v_0\times [v_+,v_{},1]`$.

The closure of $`b(X^+)`$ in $`\overline{X}`$ is $`V_0\times (V_+k)`$; denote it by $`\overline{X^+}`$, and let $`\overline{g^+}`$ be the inclusion into $`\overline{X}`$. Let $`\pi ^{}`$, $`\overline{\pi ^+}`$, $`\overline{\pi }`$ be the projections of $`X^{}`$, $`\overline{X^+}`$, and $`\overline{X}`$, respectively onto $`V_0`$.

Note that $`b(X^{})`$ is closed in $`\overline{X}`$; let $`h`$ be the inclusion of the complement. There is a distinguished triangle:

(2)
$$b_{}(g^{})_{}(g^{})^!Sb_!Sh_{}h^{}b_!S\stackrel{[1]}{}$$

Applying $`\overline{\pi }_{}`$, the left term becomes

$$\overline{\pi }_{}b_{}(g^{})_{}(g^{})^!S=(\pi ^{})_{}(g^{})^!SS^!,$$

using (1). Theorem 1 will follow if we can show that $`\overline{\pi }_{}`$ annihilates the second and third terms of (2).

For the second term we see that $`E=\overline{X}`$, $`B=\overline{X^+}`$ are of the form given in Lemma 6, so we get $`\overline{\pi }_{}b_!S=(\overline{\pi ^+})_{}(\overline{f^+})^{}b_!S`$. This vanishes, since $`S|_{X^+}=0`$ implies $`(b_!S)|_{\overline{X^+}}=0`$.

For the third term we use Lemma 6 again, this time with $`E=\overline{X}X`$, $`B=\overline{X^+}X^+`$. This is a vector bundle with only negative weights, so using the special case of Lemma 6 (with the opposite $`T`$-action), we get $`\overline{\pi }_{}h_{}h^{}b_!S\overline{\pi }_{}a^{}b_!S=0`$, where $`a`$ is the inclusion of $`\overline{X^+}X^+`$ into $`\overline{X}`$.

## 5. Proof of Theorem 2

We recall the properties of mixed Hodge modules and mixed $`l`$-adic sheaves that we will need to prove Theorem 2; more information can be found in \[BBD\], \[S\].

If $`k=`$, we let $`𝒟_m(X)`$ be Saito’s category of mixed Hodge modules. If $`k=\overline{𝔽_q}`$ and $`X`$ arises by extension of scalars from a variety $`X_0`$ defined over $`𝔽_q`$, we let $`𝒟_m(X)`$ be the category $`D_{mixed}^b(X_0;\overline{_l})`$ of mixed $`\overline{_l}`$ sheaves on $`X_0`$. These two situations share the following formal properties.

There is a triangulated functor $`r:𝒟_m(X)𝒟(X)`$ with the property that $`r(S)=0`$ if and only if $`S=0`$; it follows that any morphism $`\alpha `$ in $`𝒟_m(X)`$ is an isomorphism if and only if $`r(\alpha )`$ is.

For any $`w`$, there are full subcategories $`𝒟_w(X)`$ and $`𝒟_w(X)`$ of $`𝒟_m(X)`$. Objects of $`𝒟_w(X)𝒟_w(X)`$ are called pure of weight $`w`$. If $`Y`$ is a subvariety of $`X`$ and $`w`$, the intersection cohomology sheaf $`\mathrm{𝐈𝐂}^{}(Y)𝒟(X)`$ with constant coefficients extended by $`0`$ to $`X`$ is isomorphic to $`r(L)`$ for an object $`L𝒟_m(X)`$ which is pure of weight $`w`$. Conversely, if $`S`$ is pure, then $`r(S)`$ is a direct sum of intersection cohomology sheaves with possibly non-constant coefficients.

Given an algebraic map $`f:XY`$, the functors $`f^{}`$, $`f^!`$, $`f_{}`$, $`f_!`$ lift to functors between $`𝒟_m(X)`$ and $`𝒟_m(Y)`$ which satisfy the same adjunction properties as their non-mixed versions. We have

$$f_{}𝒟_w(X)𝒟_w(Y),f^!𝒟_w(Y)𝒟_w(X)$$

$$f_!𝒟_w(X)𝒟_w(Y),f^{}𝒟_w(Y)𝒟_w(X)$$

We say an object $`S𝒟_m(X)`$ is weakly equivariant or $`T`$-constructible if the same property holds for $`r(S)`$.

Using these properties, we can now deduce Theorem 2 from Theorem 1. The isomorphisms (1) and the transformation $`\iota _S`$ are defined by means of the standard adjunctions, so they lift to isomorphisms on weakly equivariant mixed sheaves. Thus, for any $`w`$ we have

$$𝒟_w(X)^!=(\pi ^+)_!(g^+)^{}𝒟_w(X)𝒟_w(F),\text{ and}$$

$$𝒟_w(X)^!=𝒟_w(X)^!=(\pi ^{})_{}(g^{})^!𝒟_w(X)𝒟_w(F).$$

We have proved:

###### Theorem 8.

The hyperbolic localization functors preserve purity of weakly equivariant mixed sheaves.

Theorem 2 follows.

## 6. Proof of Lemma 6

We will prove the first isomorphism of Lemma 6; the second follows by dualizing. The following proof is a generalization of an argument of Springer \[Sp\].

Applying $`\varphi _{}p_{}`$ to the adjunction map $`Si_{}i^{}S`$ gives a natural map $`\varphi _{}p_{}S\varphi _{}i^{}S`$. If it is an isomorphism for two objects in a distinguished triangle, it will also be an isomorphism for the third. If $`j`$ is the inclusion of $`Ei(B)`$ into $`E`$, we have a distinguished triangle:

$$j_!j^{}SSi_{}i^{}S\stackrel{[1]}{}$$

Since Lemma 6 is clearly true for $`i_{}i^{}S`$, it is enough to show it when $`S=j_!j^{}S`$. In this case $`i^{}S=0`$, and we must prove that $`\varphi _{}p_{}S=0`$.

Let $`\mathrm{\Gamma }_\mu `$ be the graph of the action map $`\mu :T\times EE`$, and let $`\mathrm{\Gamma }`$ be its closure in $`𝔸^1\times E\times E`$. Let $`q_1`$, $`q_2:\mathrm{\Gamma }𝔸^1\times E`$ be the restrictions of the projections $`\pi _{12}`$ and $`\pi _{13}`$. An easy computation in coordinates shows that if $`\stackrel{~}{\mathrm{\Gamma }}`$ is the closure of $`\mathrm{\Gamma }_\mu `$ in $`𝔸^1\times E\times (W)`$, then $`\stackrel{~}{\mathrm{\Gamma }}\mathrm{\Gamma }_\mu \{0\}\times E\times i(B)`$. As a result, $`\stackrel{~}{\mathrm{\Gamma }}=\mathrm{\Gamma }`$, and the map $`q_1`$ is proper.

Let $`K`$ denote the coefficient ring of our sheaves. We will write $`K_Z`$ to denote the constant rank $`1`$ local system on a variety $`Z`$, placed in degree $`0`$. Define a map $`\psi =\mathrm{id}_{𝔸^1}\times (\varphi p):𝔸^1\times E𝔸^1\times Y`$. Adjunction gives a morphism

(3) $`\psi _{}(K_{𝔸^1}S)`$ $`\psi _{}q_2q_2^{}(K_{𝔸^1}S)`$
$`\psi _{}q_{1!}q_2^{}(K_{𝔸^1}S),`$
using the properness of $`q_1`$ and the fact that $`\psi q_1=\psi q_2`$. The object on the left is isomorphic to $`K_{𝔸^1}\varphi _{}p_{}S`$, by smooth base change. Since $`q_2`$ restricts to an isomorphism over $`T\times E`$, we see that $`\alpha `$ is an isomorphism on $`T\times Y`$.

Let $`\stackrel{~}{S}`$ be the right hand side of (3). We will show that $`\stackrel{~}{S}|_{\{0\}\times Y}=0`$. Since $`\stackrel{~}{S}|_{T\times Y}K_T\varphi _{}p_{}S`$, this implies that $`\stackrel{~}{S}b_!K_T\varphi _{}p_{}S`$, where $`b:T𝔸^1`$ is the inclusion.

Composing the morphism (3) with the morphism $`b_!K_T\varphi _{}p_{}SK_{𝔸^1}\varphi _{}p_{}S`$ obtained from the adjunction morphism $`b_!K_TK_{𝔸^1}`$ gives an endomorphism of $`K_{𝔸^1}\varphi _{}p_{}S`$ which is an isomorphism over $`T\times Y`$ but which restricts to $`0`$ on $`\{0\}\times Y`$. Since this is only possible if $`\stackrel{~}{S}=0`$, we must have $`\varphi _{}p_{}S=0`$, as desired.

Finally, to show that $`\stackrel{~}{S}|_{\{0\}\times Y}=0`$, let $`E_0=Ei(B)`$, and consider the maps

where $`j`$ is the inclusion, $`\widehat{ȷ}(t,x)=(t,t^1x,x)`$ and $`\pi ,\pi ^{}`$ are the projections from $`T\times E_0`$ and $`𝔸^1\times E`$ onto the second factors. The square commutes, and is in fact Cartesian, since $`(\pi ^{}q_2)^1(E_0)\mathrm{\Gamma }_\mu `$.

Let $`S_0=S|_{E_0}`$. Then $`Sj_!S_0`$, so

$$\stackrel{~}{S}\psi _{}q_{1!}q_2^{}(\pi ^{})^{}j_!S_0\psi _{}q_{1!}\widehat{ȷ}_!\pi ^{}S_0.$$

Since $`S`$, and thus $`S_0`$, is weakly equivariant, we have

$$(q_{1!}\widehat{ȷ}_!\pi ^{}S_0)|_{T\times E_0}LS_0$$

for some local system $`L`$ on $`T`$. This implies $`q_{1!}\widehat{ȷ}_!\pi ^{}S_0b_!Lj_!S_0=b_!LS`$, since the stalks of this sheaf vanish outside of $`T\times E_0`$. Applying $`\psi _{}`$ now gives

$$\stackrel{~}{S}b_!L\varphi _{}p_{}S,$$

which implies $`\stackrel{~}{S}|_{\{0\}\times Y}=0`$, as required.


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# Spacetime structure of massive gravitino

## 1 Motivation

A massive gravitino is described by $`\psi ^\mu `$. As far as its spacetime properties are concerned, it transforms as a finite dimensional non-unitary representation of the Lorentz group,

$`\psi ^\mu :\underset{\text{Spinor}\text{Sector}}{\underset{}{\left[({\displaystyle \frac{1}{2}},0)(0,{\displaystyle \frac{1}{2}})\right]}}\underset{\text{Vector}\text{Sector}}{\underset{}{({\displaystyle \frac{1}{2}},{\displaystyle \frac{1}{2}})}}.`$ (1)
The unitarily transforming physical states are built upon this structure .

We enumerate two circumstances that motivate us to take an ab initio look at this representation space.

1. For the vector sector, it has recently been called to attention that the Proca description of the $`(1/2,1/2)`$ representation space is incomplete . An ab initio construction of this sector reveals that the Stückelberg contribution to the propagator, so important for the renormalization of the gauge theories with massive bosons , is found to naturally reside in the $`(1/2,1/2)`$ representation space.
2. At the same time, the properties of the $`(1/2,1/2)`$, along with that of the $`(1/2,0)(0,1/2)`$, representation space determine the structure of $`\psi ^\mu `$. In order to impose a single-spin, i.e., spin $`3/2`$, interpretation on the latter, the lower spin-$`1/2^+`$ and spin-$`1/2^{}`$ components of $`\psi ^\mu `$ are considered as redundant, unphysical, states that are claimed to be excluded from consideration by means of the two supplementary conditions: $`\gamma _\mu \psi ^\mu (x)=0`$, and $`_\mu \psi ^\mu (x)=0`$, respectively. However, this time-honored framework was questioned by a recent empirical observation regarding the $`N`$ and $`\mathrm{\Delta }`$ resonances . The available data on high-spin resonances reveal an unexpected and systematic clustering in terms of the $`(j/2,j/2)[(1/2,0)(0,1/2)]`$ representations with $`j=1,3`$ and $`5`$ without imposition of the supplementary conditions. For the $`N`$ and $`\mathrm{\Delta }`$ resonances these results are summarized in Fig. 1. For example, in the standard theoretical framework $`N(1440)`$, $`N(1535)`$, $`\mathrm{\Delta }(1620)`$ and $`\mathrm{\Delta }(1750)`$ should have been absent. Experimental data shows them to be present at statistically significant level.111The $`N(1440)`$, $`N(1535)`$, $`\mathrm{\Delta }(1620)`$ carry four star status, while at present $`\mathrm{\Delta }(1750)`$ simply has a one star significance .

In regard to the latter of the two enumerated circumstances, we take the position that any solution of the QCD Lagrangian for particle resonances must carry well-defined transformation properties when looked upon from different inertial frames. This forces these resonances to belong to one, or the other, of various representation spaces of the Lorentz group. For this reason the data on particle resonances may furnish hints on physical interpretation of various Lorentz group representations that one needs in gauge theories, or theories of supergravity.

For exploring the spacetime structure of massive gravitinos the charge conjugation properties play an important role. Under the operation of charge conjugation, one may choose the spinor sector to behave as a Dirac object, and implement the Majorana nature of the massive gravitino at the level of the Fock space. This is standard, see, e.g., Ref. . Or, from the very beginning choose the spinor sector to behave as a Majorana object. Since we wish to stress certain non-trivial aspects of massive gravitino that do not – at least qualitatively – depend on this choice, we shall here treat the spinor sector to be of Dirac type.

Very nature of our ab initio look at the representation space defined in Eq. (1), obliges us to present sufficient pedagogic details so that by the end of the paper much that is needed to form an opinion on the arrived results is readily available. At the same time, Letter nature of this manuscript would prevent us from delving into subtle details which are, for present, of secondary importance (but have been studied and are planned to be presented elsewhere).

We shall work in the momentum space. The notation will be essentially that introduced in Ref. .

## 2 Construction of the Spinor and Vector Sectors

We now wish to construct the primitive objects that span the representation space defined in Eq. (1) for an arbitrary $`\stackrel{}{p}`$. We first construct the spinor and vector sectors in the rest frame, and then boost them using the following boosts:

$`\kappa ^{(\frac{1}{2},0)(0,\frac{1}{2})}=\kappa ^{(\frac{1}{2},0)}\kappa ^{(0,\frac{1}{2})},\kappa ^{(\frac{1}{2},\frac{1}{2})}=\kappa ^{(\frac{1}{2},0)}\kappa ^{(0,\frac{1}{2})},`$ (2)
with

$`\kappa ^{(\frac{1}{2},0)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2m(E+m)}}}\left[(E+m)I_2+\stackrel{}{\sigma }\stackrel{}{p}\right],`$ (3)
$`\kappa ^{(0,\frac{1}{2})}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2m(E+m)}}}\left[(E+m)I_2\stackrel{}{\sigma }\stackrel{}{p}\right].`$ (4)
where $`I_n`$ stands for an $`n\times n`$ identity matrix, while the remaining symbols carry their usual contextual meaning. We define the spin-$`1/2`$ helicity operator: $`\mathrm{\Sigma }=(\stackrel{}{\sigma }/2)\widehat{p},`$ where $`\widehat{p}=\stackrel{}{p}/|\stackrel{}{p}|`$, and $`\stackrel{}{p}=|\stackrel{}{p}|(\mathrm{sin}(\theta )\mathrm{cos}(\varphi ),\mathrm{sin}(\theta )\mathrm{sin}(\varphi ),\mathrm{cos}(\theta ))`$. Keeping full freedom in the choice of phases, its positive and negative helicity states are:

$`h^+=N\mathrm{exp}\left(i\vartheta _1\right)\left(\begin{array}{c}\mathrm{cos}(\theta /2)\mathrm{exp}(i\varphi /2)\\ \mathrm{sin}(\theta /2)\mathrm{exp}(i\varphi /2)\end{array}\right),`$ (7)
$`h^{}=N\mathrm{exp}\left(i\vartheta _2\right)\left(\begin{array}{c}\mathrm{sin}(\theta /2)\mathrm{exp}(i\varphi /2)\\ \mathrm{cos}(\theta /2)\mathrm{exp}(i\varphi /2)\end{array}\right).`$ (10)
### 2.1 $`(1/2,0)(0,1/2)`$ Representation space

In this subsection we present the kinematic structure of the $`(1/2,0)(0,1/2)`$ representation space in such a manner that its extension to the $`(1/2,1/2)`$ representation space becomes transparent. The familiarity of the equation presented shall not, we hope, mar the procedure adopted. The entire construct, in essence, relies on nothing more than the boost operators.

The rest-frame $`(1/2,0)(0,1/2)`$ spinors are then chosen to be:

$`u_{+1/2}(\stackrel{}{0})=\left(\begin{array}{c}h^+\\ h^+\end{array}\right),u_{1/2}(\stackrel{}{0})=\left(\begin{array}{c}h^{}\\ h^{}\end{array}\right),`$ (15)
$`v_{+1/2}(\stackrel{}{0})=\left(\begin{array}{c}h^+\\ h^+\end{array}\right),v_{1/2}(\stackrel{}{0})=\left(\begin{array}{c}h^{}\\ h^{}\end{array}\right).`$ (20)
The choice of the phases made in writing down these spinors has been determined by the demand of parity covariance.

The boosted spinors, $`u_{\pm 1/2}(\stackrel{}{p})`$ and $`v_{\pm 1/2}(\stackrel{}{p})`$ are obtained by applying the boost operator $`\kappa ^{(\frac{1}{2},0)(0,\frac{1}{2})}`$ to the above spinors, yielding:

$`u_{+1/2}(\stackrel{}{p})={\displaystyle \frac{N\mathrm{exp}(i\vartheta _1)}{\sqrt{2m(m+E)}}}\left(\begin{array}{c}\mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\end{array}\right),`$ (25)
$`u_{1/2}(\stackrel{}{p})={\displaystyle \frac{N\mathrm{exp}(i\vartheta _2)}{\sqrt{2m(m+E)}}}\left(\begin{array}{c}\mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\end{array}\right),`$ (30)
$`v_{+1/2}(\stackrel{}{p})={\displaystyle \frac{N\mathrm{exp}(i\vartheta _1)}{\sqrt{2m(m+E)}}}\left(\begin{array}{c}\mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\end{array}\right),`$ (35)
$`v_{1/2}(\stackrel{}{p})={\displaystyle \frac{N\mathrm{exp}(i\vartheta _2)}{\sqrt{2m(m+E)}}}\left(\begin{array}{c}\mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{sin}(\theta /2)\\ \mathrm{exp}(i\varphi /2)(m+|\stackrel{}{p}|+E)\mathrm{cos}(\theta /2)\end{array}\right).`$ (40)
In the standard notation, these satisfy the orthonormality and completeness relations, along with the Dirac equation:

$`\overline{u}_h(\stackrel{}{p})u_h^{}(\stackrel{}{p})=+2N^2\delta _{hh^{}},\overline{v}_h(\stackrel{}{p})v_h^{}(\stackrel{}{p})=2N^2\delta _{hh^{}},`$ (41)
$`{\displaystyle \frac{1}{2N^2}}\left[{\displaystyle \underset{h=\pm 1/2}{}}u_h(\stackrel{}{p})\overline{u}_h(\stackrel{}{p}){\displaystyle \underset{h=\pm 1/2}{}}v_h(\stackrel{}{p})\overline{v}_h(\stackrel{}{p})\right]=I_4.`$ (42)
The wave equation satisfied by the $`u_h(\stackrel{}{p})`$ and $`v_h(\stackrel{}{p})`$ spinors follows if we note:222This part of the calculation is best done if one starts with the momenta in the Cartesian co-ordinates, and takes the “axis of spin quantization” to be the z-axis.

$`{\displaystyle \frac{1}{2N^2}}\left[{\displaystyle \underset{h=\pm 1/2}{}}u_h(\stackrel{}{p})\overline{u}_h(\stackrel{}{p})+{\displaystyle \underset{h=\pm 1/2}{}}v_h(\stackrel{}{p})\overline{v}_h(\stackrel{}{p})\right]{\displaystyle \frac{\gamma _\mu p^\mu }{m}},`$ (43)
where we defined,

$`\gamma _0=\left(\begin{array}{cc}0_2& I_2\\ I_2& 0_2\end{array}\right),\gamma _i=\left(\begin{array}{cc}0_2& \sigma _i\\ \sigma _i& 0_2\end{array}\right).`$ (48)
The $`0_n`$ are $`n\times n`$ null matrices. Adding/subtracting Eqs. (42) and (43), yields:

$`{\displaystyle \frac{1}{2N^2}}{\displaystyle \underset{h=\pm 1/2}{}}u_h(\stackrel{}{p})\overline{u}_h(\stackrel{}{p})={\displaystyle \frac{\gamma _\mu p^\mu +mI_4}{2m}},`$ (49)
$`{\displaystyle \frac{1}{2N^2}}{\displaystyle \underset{h=\pm 1/2}{}}v_h(\stackrel{}{p})\overline{v}_h(\stackrel{}{p})={\displaystyle \frac{\gamma _\mu p^\mu mI_4}{2m}}.`$ (50)
Multiplying Eq. (49) from the right by $`u_h^{}(\stackrel{}{p})`$, and Eq. (50) by $`v_h^{}(\stackrel{}{p})`$, and using Eqs. (41), immediately yield the momentum-space wave equation for the $`(1/2,0)(0,1/2)`$ representation space,

$`\left(\gamma _\mu p^\mu \pm mI_4\right)\psi _h(\stackrel{}{p})=0.`$ (51)
In Eq. (51), the minus sign is to be taken for, $`\psi _h(\stackrel{}{p})=u_h(\stackrel{}{p})`$, and the plus sign for, $`\psi _h(\stackrel{}{p})=v_h(\stackrel{}{p})`$.

The essential element to note in regard to Eq. (51) is that it follows directly from the explicit expressions for the $`u_h(\stackrel{}{p})`$ and $`v_h(\stackrel{}{p})`$.

### 2.2 $`(1/2,1/2)`$ Representation space – An ab initio construct

Next, we introduce the rest-frame vectors for the $`(1/2,1/2)`$ representation space,

$`\xi _1(\stackrel{}{0})=h^+h^+,`$ (52)
$`\xi _2(\stackrel{}{0})={\displaystyle \frac{1}{\sqrt{2}}}\left(h^+h^{}+h^{}h^+\right),`$ (53)
$`\xi _3(\stackrel{}{0})=h^{}h^{},`$ (54)
$`\xi _4(\stackrel{}{0})={\displaystyle \frac{1}{\sqrt{2}}}\left(h^+h^{}h^{}h^+\right).`$ (55)
The boosted vectors are thus: $`\xi _\zeta (\stackrel{}{p})=\kappa ^{(\frac{1}{2},\frac{1}{2})}\xi _\zeta (\stackrel{}{0}),\zeta =1,2,3,4,`$

$`\xi _1(\stackrel{}{p})={\displaystyle \frac{\mathrm{exp}(i2\vartheta _1)N^2}{2}}\left(\begin{array}{c}2\mathrm{exp}(i\varphi )\mathrm{cos}^2(\theta /2)\\ \mathrm{sin}(\theta )\\ \mathrm{sin}(\theta )\\ 2\mathrm{exp}(i\varphi )\mathrm{sin}^2(\theta /2)\end{array}\right),`$ (60)
$`\xi _2(\stackrel{}{p})={\displaystyle \frac{\mathrm{exp}(i2\vartheta _1)N^2}{\sqrt{2}m}}\left(\begin{array}{c}\mathrm{exp}(i\varphi )E\mathrm{sin}(\theta )\\ \left(|\stackrel{}{p}|+E\mathrm{cos}(\theta )\right)\\ |\stackrel{}{p}|E\mathrm{cos}(\theta )\\ \mathrm{exp}(i\varphi )E\mathrm{sin}(\theta )\end{array}\right),`$ (65)
$`\xi _3(\stackrel{}{p})={\displaystyle \frac{\mathrm{exp}(i2\vartheta _1)N^2}{2}}\left(\begin{array}{c}2\mathrm{exp}(i\varphi )\mathrm{sin}^2(\theta /2)\\ \mathrm{sin}(\theta )\\ \mathrm{sin}(\theta )\\ 2\mathrm{exp}(i\varphi )\mathrm{cos}^2(\theta /2)\end{array}\right),`$ (70)
$`\xi _4(\stackrel{}{p})={\displaystyle \frac{\mathrm{exp}(i2\vartheta _1)N^2}{\sqrt{2}m}}\left(\begin{array}{c}\mathrm{exp}(i\varphi )|\stackrel{}{p}|\mathrm{sin}(\theta )\\ \left(E+|\stackrel{}{p}|\mathrm{cos}(\theta )\right)\\ E|\stackrel{}{p}|\mathrm{cos}(\theta )\\ \mathrm{exp}(i\varphi )|\stackrel{}{p}|\mathrm{sin}(\theta )\end{array}\right).`$ (75)
In the notation of Ref. ,333In the cited work we presented the $`(1/2,1/2)`$ representation space in its parity realization. Here, the presentation is in terms of helicity realization. The two descriptions have mathematically similar but physically distinct structures, which, e.g., show up in their different behavior under the operation of Parity. these satisfy the orthonormality and completeness relations, along with a new wave equation. The orthonormality and completeness relations are:

$`\overline{\xi }_\zeta (\stackrel{}{p})\xi _\zeta ^{}(\stackrel{}{p})=N^4\delta _{\zeta \zeta ^{}},\zeta =1,2,3,`$
$`\overline{\xi }_\zeta (\stackrel{}{p})\xi _\zeta ^{}(\stackrel{}{p})=+N^4\delta _{\zeta \zeta ^{}},\zeta =4,`$
$`{\displaystyle \frac{1}{N^4}}\left[\xi _4(\stackrel{}{p})\overline{\xi }_4(\stackrel{}{p}){\displaystyle \underset{\zeta =1,2,3}{}}\xi _\zeta (\stackrel{}{p})\overline{\xi }_\zeta (\stackrel{}{p})\right]=I_4,`$ (76)
where

$`\overline{\xi }_\zeta (\stackrel{}{p})\xi _\zeta (\stackrel{}{p})^{}\lambda _{00},`$ (77)
with

$`\lambda _{00}=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 1\end{array}\right).`$ (82)
The parity operator for the $`(1/2,1/2)`$ representation space is:

$`𝒫=\lambda _{00}\mathrm{exp}[i\alpha ],`$ $`:\{\theta \pi \theta ,\varphi \pi +\varphi \}`$ (83)
$`\alpha =\text{a real number},`$
while the helicity operator for this space is, $`\stackrel{}{J}\widehat{p}`$, with $`\stackrel{}{J}`$ given by:

$`J_x={\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}0& 1& 1& 0\\ 1& 0& 0& 1\\ 1& 0& 0& 1\\ 0& 1& 1& 0\end{array}\right),J_y={\displaystyle \frac{1}{2}}\left(\begin{array}{cccc}0& i& i& 0\\ i& 0& 0& i\\ i& 0& 0& i\\ 0& i& i& 0\end{array}\right),J_z=\left(\begin{array}{cccc}1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\end{array}\right).`$ (96)
We now must take a small definitional detour towards the notion of the dragged Casimirs for spacetime symmetries. It arises in the following fashion. The second Casimir operator, $`C_2`$, of the Poincaré group is defined as the square of the Pauli-Lubanski pseudovector:

$`𝒲^\mu ={\displaystyle \frac{1}{2}}ϵ^{\mu \nu \rho \sigma }M_{\nu \rho }P_\sigma ,`$ (97)
where $`ϵ^{\mu \nu \rho \sigma }`$ is the standard Levi-Civita symbol in four dimensions, while $`M_{\mu \nu }`$ denote generators of the Lorentz group,

$`M_{0i}=K_i,M_{ij}=ϵ_{ijk}J^k,`$ (98)
where each of the $`i,j,k`$ runs over $`1,2,3`$. The $`P_\mu `$ are generators of the spacetime translations. In general, these have non-vanishing commutators with $`M_{\mu \nu }`$,

$`[P_\mu ,M_{\rho \sigma }]=i\left(\eta _{\mu \rho }P_\sigma \eta _{\mu \sigma }P_\rho \right).`$ (99)
On using Eq. (99), we rewrite $`C_2`$ as

$`C_2`$ $`=`$ $`{\displaystyle \frac{1}{4}}ϵ^{\mu \nu \rho \sigma }ϵ_{\mu \lambda \kappa \zeta }M_{\nu \rho }M^{\lambda \kappa }P_\sigma P^\zeta `$ (100)
$`+\left[{\displaystyle \frac{i}{4}}ϵ^{\mu \nu \rho \sigma }ϵ_{\mu \lambda \kappa \zeta }M_{\nu \rho }\eta _{\sigma }^{}{}_{}{}^{\lambda }P^\kappa P^\zeta +{\displaystyle \frac{i}{4}}ϵ^{\mu \nu \rho \sigma }ϵ_{\mu \kappa \lambda \zeta }M_{\nu \rho }\eta _{\sigma }^{}{}_{}{}^{\kappa }P^\lambda P^\zeta \right].`$
The squared brackets vanishes due to antisymmetry of the Levi-Civita symbol. As such, the space-time translation operators entering the definition of $`C_2`$ can be moved to the very right.

This observation allows for introducing the dragged Casimir $`\stackrel{~}{C}_2`$ as an operator with the same form as $`C_2`$ – the difference being that the commutator in Eq. (99) is now set to zero \[as is appropriate for finite dimensional $`SU_R(2)SU_L(2)`$ representations\]. Consequently, while $`C_2`$ and $`\stackrel{~}{C}_2`$ carry same invariant eigenvalues when acting upon momentum eigenstates, their commutators with the Lorentz group generators are no longer identical.444To avoid confusion, note that $`\stackrel{~}{C}_2`$ is defined in the $`SU_R(2)SU_L(2)`$; while $`C_2`$ is defined in the Poincaré group. For the $`(1/2,0)(0,1/2)`$ representation space, $`[\stackrel{~}{C_2},\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}]`$ vanishes. For the $`(1/2,1/2)`$ representation space, $`[\stackrel{~}{C_2},\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}]`$ does not vanish (except when acting upon rest states), and equals $`\mathrm{\hspace{0.17em}4}E\stackrel{}{P}\stackrel{}{K}`$. This leads to the fact that while the former representation space is endowed with a well-defined spin, the latter is not:

> As an immediate application, $`\stackrel{~}{C}_2`$ for the $`(1/2,1/2)`$ representation space bifurcates this space into two sectors. The three states $`\xi _\zeta (\stackrel{}{p})`$ with $`\zeta =1,2,3`$ are associated with the $`\stackrel{~}{C}_2`$ eigenvalue, $`\mathrm{\hspace{0.17em}2}m^2`$; while the, $`\zeta =4`$, corresponds to eigenvalue zero.
>
> Thus, all the $`\xi _\zeta (\stackrel{}{p})`$, except for the rest frame, cease to be eigenstates of the $`(1/2,1/2)`$’s $`\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}`$ and do not carry definite spins. This contrasts with the situation for the $`(1/2,0)(0,1/2)`$ representation space, where the $`\psi _h(\stackrel{}{p})`$ are eigenstates of the corresponding $`\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}`$.

Now in order that the $`\xi _\zeta (\stackrel{}{p})`$ carry the standard contravariant Lorentz index, we introduce a rotation in the $`(1/2,1/2)`$ representation space via :

$`S={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cccc}0& i& i& 0\\ i& 0& 0& i\\ 1& 0& 0& 1\\ 0& i& i& 0\end{array}\right)`$ (105)
Then, the $`(1/2,1/2)`$ representation space is spanned by four Lorentz vectors:

$`𝒜_\zeta ^\mu (\stackrel{}{p})=S^{\mu \alpha }\left[\xi _\zeta (\stackrel{}{p})\right]_\alpha ,\zeta =1,2,3,4,`$ (106)
and the superscript $`\mu `$ is the standard Lorentz index. Following the procedure established in Sec. 2.1, they can be shown to satisfy a new wave equation ,

$`\left(\mathrm{\Lambda }_{\mu \nu }p^\mu p^\nu \pm m^2I_4\right)𝒜_\zeta (\stackrel{}{p})=0,`$ (107)
where the plus sign is to be taken for, $`\zeta =1,2,3`$, while the minus sign belongs to, $`\zeta =4`$. The $`\mathrm{\Lambda }_{\mu \nu }`$ matrices are: $`\mathrm{\Lambda }_{00}=\text{diag}(1,1,1,1)`$, $`\mathrm{\Lambda }_{11}=\text{diag}(1,1,1,1)`$, $`\mathrm{\Lambda }_{22}=\text{diag}(1,1,1,1)`$, $`\mathrm{\Lambda }_{33}=\text{diag}(1,1,1,1)`$, and

$`\mathrm{\Lambda }_{01}=\left(\begin{array}{cccc}0& 1& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\mathrm{\Lambda }_{02}=\left(\begin{array}{cccc}0& 0& 1& 0\\ 0& 0& 0& 0\\ 1& 0& 0& 0\\ 0& 0& 0& 0\end{array}\right),\mathrm{\Lambda }_{03}=\left(\begin{array}{cccc}0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 0& 0& 0\\ 1& 0& 0& 0\end{array}\right),`$ (120)
$`\mathrm{\Lambda }_{12}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 1& 0\\ 0& 1& 0& 0\\ 0& 0& 0& 0\end{array}\right),\mathrm{\Lambda }_{13}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 0& 0\\ 0& 1& 0& 0\end{array}\right),\mathrm{\Lambda }_{23}=\left(\begin{array}{cccc}0& 0& 0& 0\\ 0& 0& 0& 0\\ 0& 0& 0& 1\\ 0& 0& 1& 0\end{array}\right).`$ (133)
The remaining $`\mathrm{\Lambda }_{\mu \nu }`$ are obtained from the above expressions by noting: $`\mathrm{\Lambda }_{\mu \nu }=\mathrm{\Lambda }_{\nu \mu }`$. Parenthetically, we note that the S-transformed $`\lambda _{00}`$ equals $`\mathrm{\Lambda }_{00}`$ and is nothing but the standard spacetime metric (for flat spacetime).

It can also be seen that $`\xi _\zeta (\stackrel{}{p})`$, for $`\zeta =1,2,3`$, coincide with the solutions of Proca framework (and are divergence-less); whereas $`\xi _4(\stackrel{}{p})`$, that gives the Stückelberg contribution to the propagator, lies outside the Proca framework:555We define the dragged Pauli-Lubanski pseudovector, $`\stackrel{~}{𝒲_\mu }`$, in a manner parallel to the introduction of the dragged second Casimir operator.

| $`\zeta `$ | $`p_\mu 𝒜^\mu (\stackrel{}{p})`$ | $`\stackrel{~}{𝒲_\mu }^{(1/2,1/2)}𝒜^\mu (\stackrel{}{p})`$ | $`\lambda _c`$ | Remarks |
| --- | --- | --- | --- | --- |
| $`1,2,3`$ | $`=0`$ | $`0`$ | $`2`$ | Proca Sector |
| $`4`$ | $`0`$ | $`=0`$ | $`0`$ | Stückelberg Sector |

In the above table we have introduced the $`\lambda _c`$ via the equation:

$`\stackrel{~}{C}_2^{(1/2,1/2)}𝒜(\stackrel{}{p})=\lambda _cm^2𝒜(\stackrel{}{p}).`$ (134)
## 3 Construction of spinor vector: The massive gravitino

We now wish to present the basis vectors for the representation space defined by Eq. (1) in a language which is widely used . This would allow the present analysis to be more readily available, and also bring out the relevant similarities and differences with the framework of Rarita and Schwinger .

In writing down the basis spinor-vectors, we will use the fact that in the $`(1/2,1/2)`$ representation space the charge conjugation is implemented by

$`𝒞^{(1/2,1/2)}:𝒜(\stackrel{}{p})\left[𝒜(\stackrel{}{p})\right]^{}.`$ (135)
In the $`(1/2,0)(0,1/2)`$ representation space the charge conjugation operator is $`𝒞^{(1/2,0)(0,1/2)}:i\gamma ^2K,`$ where $`K`$ complex conjugates the spinor to its right. Then, with the uniform choice,

$`\vartheta _1=\vartheta _2,`$ (136)
we obtain:

$`𝒞^{(1/2,0)(0,1/2)}:\{\begin{array}{c}u_{+1/2}(\stackrel{}{p})v_{1/2}(\stackrel{}{p}),u_{1/2}(\stackrel{}{p})v_{+1/2}(\stackrel{}{p}),\hfill \\ v_{+1/2}(\stackrel{}{p})u_{1/2}(\stackrel{}{p}),v_{1/2}(\stackrel{}{p})u_{+1/2}(\stackrel{}{p}).\hfill \end{array}`$ (139)
In the spirit outlined, the massive gravitino lives in a space spanned by sixteen spinor-vectors defined in items A, B, C below:

1. Of these, eight spinor-vectors have $`\stackrel{~}{C}_2`$ – but not $`\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}`$ – eigenvalues, $`\frac{15}{4}m^2`$. These can be further subdivided into particle,

$`\psi _a^\mu (\stackrel{}{p}):\{\begin{array}{cc}& \psi _1^\mu (\stackrel{}{p})=u_{+1/2}(\stackrel{}{p})𝒜_1^\mu (\stackrel{}{p})\text{ ,}\hfill \\ & \psi _2^\mu (\stackrel{}{p})=\sqrt{\frac{2}{3}}u_{+1/2}(\stackrel{}{p})𝒜_2^\mu (\stackrel{}{p})+\sqrt{\frac{1}{3}}u_{1/2}(\stackrel{}{p})𝒜_1^\mu (\stackrel{}{p})\text{ ,}\hfill \\ & \psi _3^\mu (\stackrel{}{p})=\sqrt{\frac{1}{3}}u_{+1/2}(\stackrel{}{p})𝒜_3^\mu (\stackrel{}{p})+\sqrt{\frac{2}{3}}u_{1/2}(\stackrel{}{p})𝒜_2^\mu (\stackrel{}{p})\text{ ,}\hfill \\ & \psi _4^\mu (\stackrel{}{p})=u_{1/2}(\stackrel{}{p})𝒜_3^\mu (\stackrel{}{p})\text{ ,}\hfill \end{array}`$
and antiparticle sectors:

$`[\psi _a^\mu (\stackrel{}{p})]^𝒞:\{\begin{array}{cc}& \psi _5^\mu (\stackrel{}{p})=v_{1/2}(\stackrel{}{p})\left[𝒜_1^\mu (\stackrel{}{p})\right]^{}\text{ ,}\hfill \\ & \psi _6^\mu (\stackrel{}{p})=\sqrt{\frac{2}{3}}v_{1/2}(\stackrel{}{p})\left[𝒜_2^\mu (\stackrel{}{p})\right]^{}+\sqrt{\frac{1}{3}}v_{+1/2}(\stackrel{}{p})\left[𝒜_1^\mu (\stackrel{}{p})\right]^{}\text{ ,}\hfill \\ & \psi _7^\mu (\stackrel{}{p})=\sqrt{\frac{1}{3}}v_{1/2}(\stackrel{}{p})\left[𝒜_3^\mu (\stackrel{}{p})\right]^{}+\sqrt{\frac{2}{3}}v_{+1/2}(\stackrel{}{p})\left[𝒜_2^\mu (\stackrel{}{p})\right]^{}\text{ ,}\hfill \\ & \psi _8^\mu (\stackrel{}{p})=v_{+1/2}(\stackrel{}{p})\left[𝒜_3^\mu (\stackrel{}{p})\right]^{}\text{ .}\hfill \end{array}`$
Here, $`[\psi _\tau ^\mu (\stackrel{}{p})]^𝒞=𝒞^{(1/2,0)(0,1/2)}𝒞^{(1/2,1/2)}\psi _\tau ^\mu (\stackrel{}{p})`$, $`\tau =a,b,c`$.
2. Four spinor-vectors have $`\stackrel{~}{C}_2`$ – but not $`\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}`$ – eigenvalues, $`\frac{3}{4}m^2`$:

$`\psi _b^\mu (\stackrel{}{p}):\{\begin{array}{cc}& \psi _9^\mu (\stackrel{}{p})=\sqrt{\frac{2}{3}}u_{1/2}(\stackrel{}{p})𝒜_1^\mu (\stackrel{}{p})\sqrt{\frac{1}{3}}u_{+1/2}(\stackrel{}{p})𝒜_2^\mu (\stackrel{}{p})\text{ ,}\hfill \\ & \psi _{10}^\mu (\stackrel{}{p})=\sqrt{\frac{1}{3}}u_{1/2}(\stackrel{}{p})𝒜_2^\mu (\stackrel{}{p})\sqrt{\frac{2}{3}}u_{+1/2}(\stackrel{}{p})𝒜_3^\mu (\stackrel{}{p})\text{ ,}\hfill \end{array}`$
$`[\psi _b^\mu (\stackrel{}{p})]^𝒞:\{\begin{array}{cc}& \psi _{11}^\mu (\stackrel{}{p})=\sqrt{\frac{2}{3}}v_{+1/2}(\stackrel{}{p})\left[𝒜_1^\mu (\stackrel{}{p})\right]^{}+\sqrt{\frac{1}{3}}v_{1/2}(\stackrel{}{p})\left[𝒜_2^\mu (\stackrel{}{p})\right]^{}\text{ ,}\hfill \\ & \psi _{12}^\mu (\stackrel{}{p})=\sqrt{\frac{1}{3}}v_{+1/2}(\stackrel{}{p})\left[𝒜_2^\mu (\stackrel{}{p})\right]^{}+\sqrt{\frac{2}{3}}v_{1/2}(\stackrel{}{p})\left[𝒜_3^\mu (\stackrel{}{p})\right]^{}\text{ .}\hfill \end{array}`$
3. Another set of four spinor-vectors with $`\stackrel{~}{C}_2`$ – but not $`\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}`$ – eigenvalues, $`\frac{3}{4}m^2`$:

$`\psi _c^\mu (\stackrel{}{p}):\{\begin{array}{cc}& \psi _{13}^\mu (\stackrel{}{p})=u_{+1/2}(\stackrel{}{p})𝒜_4^\mu (\stackrel{}{p})\text{ ,}\hfill \\ & \psi _{14}^\mu (\stackrel{}{p})=u_{1/2}(\stackrel{}{p})𝒜_4^\mu (\stackrel{}{p})\text{ ,}\hfill \end{array}`$
$`[\psi _c^\mu (\stackrel{}{p})]^𝒞:\{\begin{array}{cc}& \psi _{15}^\mu (\stackrel{}{p})=v_{1/2}(\stackrel{}{p})\left[𝒜_4^\mu (\stackrel{}{p})\right]^{}\text{ ,}\hfill \\ & \psi _{16}^\mu (\stackrel{}{p})=v_{+1/2}(\stackrel{}{p})\left[𝒜_4^\mu (\stackrel{}{p})\right]^{}\text{ .}\hfill \end{array}`$

We have evaluated $`\gamma _\mu \psi ^\mu (\stackrel{}{p})`$, $`p_\mu \psi ^\mu (\stackrel{}{p})`$, and $`\stackrel{~}{𝒲}_\mu ^{(1/2,1/2)}\psi ^\mu (\stackrel{}{p})`$, for all of the above sixteen spinor vectors. The $`p_\mu \psi ^\mu (\stackrel{}{p})`$, when transformed to the configuration space, tests the divergence of $`\psi ^\mu (x)`$.

For $`\eta =1,4,5,8`$, $`\gamma _\mu \psi ^\mu (\stackrel{}{p})`$ identically vanishes. Requiring it to vanish for $`\eta =2,3,6,7`$ results in:

1. $`\vartheta _1=\vartheta _2.`$ When combined with Eq. (136) this implies, $`\vartheta _1=0=\vartheta _2`$. As a consequence, the global phase factors $`\vartheta _1`$ and $`\vartheta _2`$ that appear in Eqs. (10), (40), and (75) are not entirely free.
2. $`E^2=|\stackrel{}{p}|^2+m^2`$

The $`\tau =b,c`$ sectors, if (wrongly) imposed with the vanishing of, $`\gamma _\mu \psi ^\mu (\stackrel{}{p})`$ and $`p_\mu \psi ^\mu (\stackrel{}{p})`$, results in kinematically acausal dispersion relation (i.e., in $`E^2|\stackrel{}{p}|^2+m^2`$). This could be the source of the well-known problems of the Rarita-Schwinger framework as noted in works of Johnson and Sudarshan , and those of Velo and Zwanziger . In this context one may wish to recall that interactions can induce transitions between different $`\tau `$ sectors.

The analysis for all the $`\tau `$ sectors of the $`\psi ^\mu (\stackrel{}{p})`$ can be summarized in the following table:

| $`\tau `$ | $`p_\mu \psi ^\mu (\stackrel{}{p})`$ | $`\gamma _\mu \psi ^\mu (\stackrel{}{p})`$ | $`\stackrel{~}{𝒲}_\mu ^{(1/2,1/2)}\psi ^\mu (\stackrel{}{p})`$ | $`\lambda _c`$ | Remarks |
| --- | --- | --- | --- | --- | --- |
| $`a`$ | $`=0`$ | $`=0`$ | $`0`$ | $`\frac{15}{4}`$ | Rarita-Schwinger Sector |
| $`b`$ | $`=0`$ | $`0`$ | $`0`$ | $`\frac{3}{4}`$ |  |
| $`c`$ | $`0`$ | $`0`$ | $`=0`$ | $`\frac{3}{4}`$ |  |

The table clearly illustrates that there is no particular reason – except (the unjustified) insistence that each particle of nature be associated with a definite spin – to favor one $`\tau `$ sector over the other. Each of the $`\tau `$ sectors is endowed with specific properties. The Rarita-Schwinger sector has no more, or no less, physical significance than the other two sectors. While, for instance, the Rarita-Schwinger sector can be characterized by vanishing of the $`p_\mu \psi ^\mu (\stackrel{}{p})`$ and $`\gamma _\mu \psi ^\mu (\stackrel{}{p})`$; the $`\tau =c`$ sector is uniquely characterized by vanishing of $`\stackrel{~}{𝒲}_\mu ^{(1/2,1/2)}\psi ^\mu (\stackrel{}{p})`$. The $`\tau =b`$ sector allows for vanishing of $`p_\mu \psi ^\mu (\stackrel{}{p})`$ only.

Except for the rest frame, the $`\psi _\tau ^\mu (\stackrel{}{p})`$, in general, are not eigenstates of the $`\stackrel{}{J}^{\mathrm{\hspace{0.17em}2}}`$ for representation space (1). Instead, the three $`\tau `$ sectors of the representation space under consideration correspond to the following inertial-frame independent values of the associated dragged second Casimir invariant:

$`\stackrel{~}{C}_2^{[(1/2,0)(0,1/2)](1/2,1/2)}\psi _\tau ^\mu (\stackrel{}{p})=m^2\lambda _c\psi _\tau ^\mu (\stackrel{}{p}).`$ (140)
For each of the $`\tau `$ sectors, the $`\lambda _c`$ are given in the table above. Stated differently, the $`\tau =b,c`$ sectors do not carry spin one half. Similarly, the, $`\tau =a`$, sector is not a spin three half sector. The consequence is that the $`\tau =b,c`$ sector, in particular, should not be treated as a Dirac representation space. The correct wave equation for $`\psi ^\mu (\stackrel{}{p})`$ is:

$`\left[\left(\mathrm{\Lambda }_{\mu \nu }p^\mu p^\nu \pm m^2I_4\right)\left(\gamma _\mu p^\mu \pm mI_4\right)\right]\psi ^ϵ(\stackrel{}{p})=0.`$ (141)
In the standard Rarita-Schwinger framework $`_\mu \psi ^\mu (x)`$ and $`\gamma _5\gamma _\mu \psi ^\mu (x)`$ do indeed behave as Dirac spinors, and do indeed satisfy the Dirac equation. However, they are not identical to the $`\tau =b,c`$ sectors (which do not carry a characterization in terms of spin one half). If one (mistakenly) makes this identification, and sets $`_\mu \psi ^\mu (x)`$ and $`\gamma _5\gamma _\mu \psi ^\mu (x)`$ to zero, one introduces an element of kinematic acausality.

## 4 Interpretation

If one is to respect the mathematical completeness of the spinor-vector representation space associated with $`\psi ^\mu (x)`$, the Rarita-Schwinger framework cannot be considered to describe the full physical content of the representation space associated with a massive gravitino. This circumstance is akin to Dirac’s observation that a part of a representation space \[which would have violated the completeness of the $`(1/2,0)(0,1/2)`$\] cannot be “projected out” without introducing certain mathematical inconsistencies, and loosing its physical content (i.e. antiparticle, or particles).666Similar remarks apply to Majorana construct in the $`(1/2,0)(0,1/2)`$ representation . Further, the same qualitative remarks apply to the $`(1/2,1/2)`$ representation space when in the Proca framework one only confines to the divergence-less vectors. The “projecting out” of the divergence-full vector, throws away the Stückelberg contribution to the propagator, and in addition leaves the $`(1/2,1/2)`$ representation space mathematically incomplete. Now, we suggest that for the representation space defined by Eq. (1), one needs to consider all three $`\tau `$ sectors of $`\psi ^\mu (x)`$ as physical, and necessary for its mathematical consistency. The suggested framework already carries consistency with the known data on the $`N`$ and $`\mathrm{\Delta }`$ resonances, and asks that massive gravitino be considered as an object that is better described by the eigenvalues of the dragged second Casimir operator. In its rest frame it is endowed with a spin three half, and two spin half, components. A spin measurement for unpolarized ensemble of massive gravitinos at rest would yield the results $`3/2`$ with probability one half, and $`1/2`$ with probability one half. The latter probability is distributed uniformly, i.e. as one quarter, over each of the, $`\tau =b`$, and, $`\tau =c`$, sectors.


text:
# 1 Introduction

## 1 Introduction

The notion of Poincaré recurrence in classical mechanics is quite well-known. Roughly it means that within experimental error a classical system confined to a finite volume in phase space will eventually return to its initial state. This happens because of Liouville’s theorem which states that Lebesgue measure is invariant under the Hamiltonian flow.

Recurrence also occurs in quantum mechanics. One approach to recurrence in quantum mechanics has been through the theory of almost periodic functions (see for example ,  and ). Another line of research, involving coherent states, along with possible applications of quantum recurrence, can be traced in  and references therein. However, these methods differ considerably from the measure theoretic techniques employed to study recurrence in classical mechanics.

In this paper we intend to show how recurrence in quantum mechanics can be cast in a mathematical form that looks the same as the classical case. More precisely, the quantum case is a non-commutative extension of the classical case. Some of the methods presented also provide a general view on how to translate between the quantum and classical descriptions of nature.

A few remarks concerning the mathematical setting are in order. Recently C.P. Niculescu, A. Ströh and L. Zsidó in , working from a purely mathematical viewpoint, showed that an analogue of Poincaré recurrence can be obtained in a C\*-algebraic framework. Since both quantum and classical mechanics can be formulated in the language of C\*-algebras, in seems most natural to work in this setting. In fact, as we shall see in Section 2, quantum mechanics and classical mechanics are identical, except for commutativity, when both are viewed purely in C\*-algebraic terms. Our approach to Poincaré recurrence will differ somewhat from that of  in that we will also consider mappings between C\*-algebras, rather than just linear functionals on C\*-algebras. Furthermore, instead of looking at arbitrary elements of the algebras, we will concentrate on the projections. The reasons for this will become clear in Sections 2 and 3. The main mathematical results are presented in Section 4.

For these results to have implications for quantum mechanics, we can expect from our remarks concerning the classical case that we will need a quantum mechanical analogue of Liouville’s theorem. We propose such an analogue in Section 3, and in the process we are naturally led to consider finite von Neumann algebras. In Section 5 we describe how the theorems of Section 4 would result in recurrence in quantum mechanics. Using the analogy between quantum and classical mechanics we also briefly discuss the properties a quantum mechanical system should most likely have in order to satisfy the requirements of these theorems.

## 2 Quantum mechanics and classical mechanics in a C\*-algebraic setting

We start with two simple definitions that apply to both quantum mechanics and classical

mechanics:

Definition 2.1. An observable of a physical system is any attribute of the system which results in a real number when measured. We call this real number the value of the

observable during the measurement.

Definition 2.2. Consider any observable of a physical system, and any Borel set $`S`$. We now perform an experiment on the system which results in a “yes” if the value of the observable lies in $`S`$ during the experiment, and a “no” otherwise; the experiment gives no further information. We call this a yes/no experiment.

Definition 2.2 seems justified since in practice there are always experimental errors, in other words we always get a range of values (namely $`S`$ in Definition 2.2) rather than a single value.

Let’s look at the C\*-algebraic formulation of quantum mechanics (also see ). Consider any quantum mechanical system. We represent the observables of the system by a unital C\*-algebra $`𝔄`$, called the observable algebra of the system, and the state of the system by a state $`\omega `$ on $`𝔄`$ (i.e. $`\omega `$ is a normalized positive linear functional on $`𝔄`$). $`𝔄`$ contains the spectral projections of the system’s observables rather than the observables themselves. By this we mean the following: To any yes/no experiment that we can perform on the system, there corresponds a projection $`P`$ in $`𝔄`$ such that $`\omega (P)`$ is the probability of getting a “yes” during the experiment for any state $`\omega `$ of the system. We will refer to $`P`$ as the projection of the yes/no experiment.

We will only consider yes/no experiments for which the experimental setup is such that at least in the case of a “yes” the system survives the experiment (for example, it is not absorbed by a detector), so further experiments can be performed on it. What does the system’s state look like after such an experiment? Consider for the moment the Hilbert space setting for quantum mechanics. Here the (pure) states of a system are represented by non-zero vectors in a Hilbert space $``$, called the state space of the system. Suppose the state is given by the unit vector $`x`$ in $``$. After a yes/no experiment the state is given by the projection of $`x`$ on some Hilbert subspace of $``$. Denoting the corresponding projection operator in case of a “yes” by $`Q`$, we see that the system’s state after the experiment would then be given by the unit vector $`Qx/Qx`$. It is clear that $`Q`$ is the projection of the experiment, since $`Qx^2=x,Qx`$ is exactly the probability of getting a “yes”. (Here the state $`\theta `$ on the C\*-algebra $`𝔏()`$ of all bounded linear operators on $``$, given by $`\theta (A)=x,Ax`$, is the C\*-algebraic representation of the state $`x`$, in the sense of $`\omega `$ above.)

Returning to our system with observable algebra $`𝔄`$, we know by the GNS-construction (see for example Section 2.3.3 of ) that there exists a Hilbert space $``$, a $``$-homomorphism $`\pi :𝔄𝔏()`$, and a unit vector $`\mathrm{\Omega }`$ in $``$, such that

$$\omega (A)=\mathrm{\Omega },\pi (A)\mathrm{\Omega }$$
(1)
for all $`A`$ in $`𝔄`$. This looks like the usual expression for the expectation value of an observable (here represented by $`\pi (A)`$) for a system in the state $`\mathrm{\Omega }`$ in the Hilbert space setting (compare $`\theta `$ above). On a heuristic level we therefore regard $``$ as the state space of the system, and $`\mathrm{\Omega }`$ as its state. Say the result of the yes/no experiment with projection $`P`$ is “yes”. On the basis of the Hilbert space setting described above, it would now be natural to expect that after the experiment the state is represented by the unit vector $`\mathrm{\Omega }^{}=\pi (P)\mathrm{\Omega }/\pi (P)\mathrm{\Omega }`$, since $`\pi (P)`$ is the projection of the experiment in the Hilbert space setting in the same way as $`Q`$ above (and hence $`\pi (P)`$ here plays the role of $`Q`$). Note that $`\pi (P)\mathrm{\Omega }^2=\omega (P)>0`$ since this is exactly the probability of getting the result “yes”. We now replace $`\mathrm{\Omega }`$ in (1) by $`\mathrm{\Omega }^{}`$ to get a new expectation functional $`\omega ^{}`$ defined by

$$\omega ^{}(A)=\mathrm{\Omega }^{},\pi (A)\mathrm{\Omega }^{}$$

for all $`A`$ in $`𝔄`$. Clearly $`\omega ^{}(A)=\omega (PAP)/\omega (P)`$, so $`\omega ^{}(1)=1`$, which implies that $`\omega ^{}`$ is a state on $`𝔄`$. Based on these arguments we give the following postulate:

Postulate 2.3. Consider a quantum mechanical system in the state $`\omega `$ on its observable algebra $`𝔄`$. Suppose we get a “yes” during a yes/no experiment performed on the system. After the experiment the state of the system is then given by the state $`\omega ^{}`$ on $`𝔄`$ defined by

$$\omega ^{}(A)=\omega (PAP)/\omega (P)$$

for all $`A`$ in $`𝔄`$, where $`P`$ is the projection of the yes/no experiment.

When expressed in terms of a density operator $`\rho `$ on a Hilbert space, where $`\omega (A)=`$ Tr$`(\rho A)`$ for a bounded linear operator $`A`$ on the Hilbert space, this is sometimes referred to as the Lüders rule (see  or ).

Lastly we mention that the time-evolution of the system is described by a one-parameter $``$-automorphism group $`\tau `$ of $`𝔄`$, so if the projection of a yes/no experiment is $`P`$ at time $`0`$, then at time $`t`$ the projection of the same yes/no experiment will be $`\tau _t(P)`$.

Now we turn to classical mechanics. We can represent the state of a classical system by a point in its phase space $`^{2n}`$. This is somewhat restrictive since such a point represents exact knowledge of the state of the system, which is impossible in practice. Therefore we rather represent the state of the system by a Borel measure $`\mu `$ on $`^{2n}`$ such that $`\mu (S)`$ is the probability that the system’s state is a point somewhere in the Borel set $`S^{2n}`$. In particular we have $`\mu (^{2n})=1`$.

We view each observable of the system as a Borel function $`f:^{2n}`$. This simply means that if the system’s state is the point $`x`$ in $`^{2n}`$, then the value of the observable is $`f(x)`$. If we perform a yes/no experiment to determine if $`f`$’s value lies in the Borel set $`S`$, then the probability of getting “yes” is clearly

$$\mu \left(f^1(S)\right)=\chi _{f^1(S)}𝑑\mu $$

where $`\chi `$ denotes characteristic functions (i.e. for any set $`A`$, the function $`\chi _A`$ assumes the value $`1`$ on $`A`$, and zero everywhere else). We can view $`\chi _{f^1(S)}`$ as a spectral projection of the observable $`f`$, and we will refer to it as the projection of the yes/no experiment, just as in the quantum mechanical case. Note that $`\chi _{f^1(S)}`$ is a projection in the C\*-algebra $`B_{\mathrm{}}(^{2n})`$ of all bounded complex-valued Borel functions on $`^{2n}`$. We can define a state $`\omega `$ on the C\*-algebra $`B_{\mathrm{}}(^{2n})`$ by

$$\omega (g)=g𝑑\mu $$

for all $`g`$ in $`B_{\mathrm{}}(^{2n})`$. Then we see that the probability of getting a “yes” in the above mentioned yes/no experiment is $`\omega (\chi _{f^1(S)})`$. So we can view $`\omega `$ as representing the state of the system in exactly the same way as in quantum mechanics, where now $`B_{\mathrm{}}(^{2n})`$ is the unital C\*-algebra representing the observables of the system. For this reason we call $`B_{\mathrm{}}(^{2n})`$ the observable algebra of the system. Postulate 2.3 then holds for the classical case as well since a “yes” will mean the system’s state is a point in $`f^1(S)`$, in which case we can describe the system’s state after the experiment by the measure $`\mu ^{}`$ given by

$$\mu ^{}(V)=\mu (Vf^1(S))/\mu (f^1(S))$$

for all Borel sets $`V^{2n}`$. As in the case of $`\mu `$ and $`\omega `$ above, $`\mu ^{}`$ corresponds to the state $`\omega ^{}`$ on $`B_{\mathrm{}}(^{2n})`$ given by

$$\omega ^{}(g)=g𝑑\mu ^{}=\omega (\chi _{f^1(S)}g\chi _{f^1(S)})/\omega (\chi _{f^1(S)})$$

(the second equality follows using standard measure theoretic arguments, i.e. first prove it for $`g`$ a characteristic function and then use Lebesgue convergence). This is exactly what Postulate 2.3 says if we replace the word “quantum” by “classical”.

For the time-evolution of a classical system we need the concept of a flow. Consider a measure space $`(X,\mathrm{\Sigma },\mu )`$, where $`\mu `$ is a measure defined on a $`\sigma `$-algebra $`\mathrm{\Sigma }`$ of subsets of the set $`X`$. A flow on $`(X,\mathrm{\Sigma },\mu )`$ is a mapping $`tT_t`$ on $``$ with the following properties: $`T_t`$ is a function defined on $`X`$ to itself, $`T_0`$ is the identity on $`X`$ (i.e. $`T_0(x)=x`$), $`T_sT_t=T_{s+t}`$, and $`T_t(S)\mathrm{\Sigma }`$ and $`\mu (T_t(S))=\mu (S)`$ for all $`S`$ in $`\mathrm{\Sigma }`$. We denote this flow simply by $`T_t`$.

The time-evolution of our classical system is given by a flow $`T_t`$ on $`(^{2n},,\lambda )`$, where $``$ is the $`\sigma `$-algebra of Borel sets of $`^{2n}`$, and $`\lambda `$ is the Lebesgue measure on $`^{2n}`$. Note that this statement contains Liouville’s theorem, namely $`\lambda (T_t(S))=\lambda (S)`$ for all $`S`$ in $``$. We call $`T_t`$ the Hamiltonian flow. It simply means that if at time $`0`$ the system is in the state $`x^{2n}`$, then at time $`t`$ it is in the state $`T_t(x)`$.

As in the C\*-algebraic approach to quantum mechanics, we want the time-evolution to act on the observable algebra rather than on the states. It is clear that an observable given by $`f`$ at time $`0`$, will then be given by $`fT_t`$ at time $`t`$ (the well-known Koopman construction, ). This is equivalent to the action of $`T_t`$ on the spectral projections of $`f`$, since $`\chi _{(fT_t)^1(S)}=\chi _{f^1(S)}T_t`$ for all Borel sets $`S`$. It is easily seen that if we define $`\tau `$ by

$$\tau _t(g)=gT_t$$
(2)
for all $`g`$ in $`B_{\mathrm{}}(^{2n})`$, then $`\tau `$ is a one-parameter $``$-automorphism group of the C\*-algebra $`B_{\mathrm{}}(^{2n})`$. So the time-evolution is described in exactly the same way as in quantum mechanics when we are working in the C\*-algebraic setting.

We have now obtained a C\*-algebraic formulation of classical mechanics. Note that $`B_{\mathrm{}}(^{2n})`$ is an abelian C\*-algebra. Replacing $`B_{\mathrm{}}(^{2n})`$ by an arbitrary abelian unital C\*-algebra would give us an abstract C\*-algebraic formulation of classical mechanics. From our discussion above it is clear that if in the C\*-algebraic formulation of quantum mechanics described earlier we assume that $`𝔄`$ is abelian, then we get exactly this abstract C\*-algebraic formulation of classical mechanics. Setting $`𝔄=B_{\mathrm{}}(^{2n})`$ would make it concrete. In this sense the C\*-algebraic formulation of quantum mechanics actually contains classical mechanics as a special case.

## 3 A quantum mechanical analogue of Liouville’s theorem

We have seen in Section 2 that in purely C\*-algebraic terms quantum mechanics and classical mechanics are identical, except of course for the fact that the classical observable algebra is abelian while this is not in general true for quantum mechanics. This suggests that it might be possible to find a quantum mechanical analogue of Liouville’s theorem. Our first clue in this direction is the following simple proposition, which is proved by standard measure theoretic arguments:

Proposition 3.1. Let $`(X,\mathrm{\Sigma },\mu )`$ be a measure space with $`\mu (X)<\mathrm{}`$, and let $`T:XX`$ be a mapping such that $`T^1(S)\mathrm{\Sigma }`$ for all $`S\mathrm{\Sigma }`$. Let $`B_{\mathrm{}}(\mathrm{\Sigma })`$ be the C\*-algebra of all bounded complex-valued $`\mathrm{\Sigma }`$-measurable functions on $`X`$, and define $`\tau `$ and $`\phi `$ by $`\tau (g)=gT`$ and $`\phi (g)=g𝑑\mu `$ for all $`gB_{\mathrm{}}(\mathrm{\Sigma })`$. Then $`\mu (T^1(S))=\mu (S)`$ for all $`S\mathrm{\Sigma }`$ if and only if $`\phi (\tau (g))=\phi (g)`$ for all $`gB_{\mathrm{}}(\mathrm{\Sigma })`$.

Consider a classical system confined to a bounded Borel set $`F`$ in the phase space $`^{2n}`$. So $`\lambda (F)<\mathrm{}`$, where $`\lambda `$ is the Lebesgue measure on $`^{2n}`$. We define a measure $`\nu `$ on the Borel sets of $`^{2n}`$ by

$$\nu (S)=\lambda (SF)\text{.}$$

Using Proposition 3.1 we see that Liouville’s theorem for this system can then be expressed in C\*-algebraic terms by stating that

$$\phi (\tau _t(g))=\phi (g)$$
(3)
for all $`g`$ in $`B_{\mathrm{}}(^{2n})`$, where $`\tau `$ is given by (2), and $`\phi (g)=g𝑑\nu `$ (so $`\phi `$ is a positive linear functional on $`B_{\mathrm{}}(^{2n})`$). Note that the condition $`\mu (X)<\mathrm{}`$ in Proposition 3.1 can be dropped if we only consider positive elements of $`B_{\mathrm{}}(\mathrm{\Sigma })`$. Hence (3) would express Liouville’s theorem for systems not necessarily bounded in phase space if we were to use $`\lambda `$ instead of $`\nu `$, and only consider positive elements $`g`$ of $`B_{\mathrm{}}(^{2n})`$. (In this case $`\phi `$ could assume infinite values and it would not be a linear mapping on $`B_{\mathrm{}}(^{2n})`$ any more.) We only work with the bounded case in recurrence though.

Because of Section 2, we now suspect that a quantum mechanical analogue of Liouville’s theorem should have the same form as (3). Let’s look at this from a different angle. In the Hilbert space setting for quantum mechanics, the state space $``$ can be viewed as the analogue of the classical phase space $`^{2n}`$. $``$ is a Hilbert space while we view $`^{2n}`$ purely as a measure space. Apart from dynamics, we saw in Section 2 that the central objects in both quantum and classical mechanics are the projections. A projection defined on $``$ is equivalent to a Hilbert subspace of $``$ (namely the range of the projection). A projection defined on $`^{2n}`$ is a Borel measurable characteristic function, and is therefore equivalent to a Borel set in $`^{2n}`$. Liouville’s theorem is based on the existence of a natural way of measuring the size of a Borel set in $`^{2n}`$, namely the Lebesgue measure $`\lambda `$. We would therefore like to have a natural way of measuring the size of a Hilbert subspace of $``$ in order to get a quantum analogue of Liouville’s theorem. An obvious candidate is the (Hilbert) dimension $`dim`$. For the Hamiltonian flow $`T_t`$, Liouville’s theorem states that $`\lambda (T_t(S))=\lambda (S)`$ for every Borel set $`S`$. (We use $`T_t(S)`$ instead of $`T_t(S)`$, since this corresponds to the action of $`T_t`$ on the observable algebra rather than on the states, namely $`\chi _ST_t=\chi _{T_t(S)}`$.) In the state space time-evolution is given by a one-parameter unitary group $`U_t`$ on $``$, and for any Hilbert subspace $`𝔎`$ of $``$ we have $`dim(U_t^{}𝔎)=dim(U_t𝔎)=dim(𝔎)`$. This is clearly similar to Liouville’s theorem. For a finite dimensional state space we will in fact view this as a quantum analogue of Liouville’s theorem. However, since state spaces are usually infinite dimensional, we would like to work with something similar to $`dim`$ which does not assume infinite values.

This leads us naturally to the C\*-algebras known as finite von Neumann algebras (see for example ), since for such an algebra there is a dimension function , defined on the projections of the algebra, which does not assume infinite values. This function is in fact the restriction of a so-called trace defined on the whole algebra, so we might as well work with this trace. We now explain this in more detail.

Let $`𝔐`$ be a finite von Neumann algebra on a Hilbert space $``$, and let $`𝔐^{}`$ be its commutant. Then there is a unique positive linear mapping tr$`:𝔐𝔐𝔐^{}`$ such that tr$`(AB)=`$ tr$`(BA)`$ and tr$`(C)=C`$ for all $`A,B𝔐`$ and $`C𝔐𝔐^{}`$. We call tr the trace of $`𝔐`$. We mention that in the special case where $`𝔐=𝔏()`$, with $``$ finite dimensional, tr is just the usual trace (sum of eigenvalues) normalized such that tr$`(1)=1`$.

For a projection $`P𝔐`$ of $``$ onto the Hilbert subspace $`𝔎`$, we see that $`U_t^{}PU_t`$ is the projection of $``$ onto $`U_t^{}𝔎`$, where $`U_t`$ is a one-parameter unitary group on $``$. So in the framework of finite von Neumann algebras we would like to replace the equation $`dim(U_t^{}𝔎)=dim(𝔎)`$ mentioned above by tr$`(U_t^{}PU_t)=`$ tr$`(P)`$.

If a self-adjoint (possibly unbounded) operator $`A`$ in $``$ is an observable and $`𝔐`$ an observable algebra of a physical system, then we want the spectral projections $`\chi _S(A)`$ of $`A`$ to be contained in $`𝔐`$, where $`S`$ is any Borel set in $``$, since these projections are the projections of the yes/no experiments that can be performed on the system. But then $`f(A)𝔐`$ for any bounded complex-valued Borel function $`f`$on $``$. In particular $`e^{iAt}𝔐`$ for all real $`t`$.

For these reasons we will consider physical systems of the following nature:

Definition 3.2. A bounded quantum system is a quantum mechanical system for which we can take the observable algebra as a finite von Neumann algebra $`𝔐`$ on a Hilbert space $``$ such that the Hamiltonian $`H`$ of the system is a self-adjoint (possibly unbounded) operator in $``$ with $`e^{iHt}𝔐`$ for real $`t`$. We denote this system by $`(𝔐,,H)`$.

The reason for the term “bounded” will become clear in Section 5. We now propose a quantum analogue of Liouville’s theorem based on the intuitive arguments in terms of dimension given above. We give it in the form of a proposition:

Proposition 3.3. Consider a bounded quantum system $`(𝔐,,H)`$. Then $`U_t=e^{iHt}`$ is a one-parameter unitary group on $``$. Let $`\tau `$ be the time-evolution of the system, i.e. $`\tau _t(A)=U_t^{}AU_t`$ for all $`A𝔐`$. Then

$$\text{tr}(\tau _t(A))=\text{ tr}(A)$$
(4)
for all $`A`$ in $`𝔐`$, where tr is the trace of $`𝔐`$. (This last statement is our quantum analogue of Liouville’s theorem.)

Proof. Since $`U_t𝔐`$, we have tr$`(\tau _t(A))=`$ tr$`(U_t^{}AU_t)=`$ tr$`(U_tU_t^{}A)=`$ tr$`(A)`$. $`\mathrm{}`$

As we suspected, our quantum analogue of Liouville’s theorem, expressed by (4), is of the same form as the C\*-algebraic formulation of the classical Liouville theorem as given by (3), with $`\phi `$ replaced by tr. Remember that $`\phi `$ and tr are both positive linear mappings on the respective observable algebras.

Remark. The classical Liouville theorem can also be expressed in terms of the Liouville equation

$$\frac{\rho }{t}=\{\rho ,H\}$$

where $`\rho :^{2n}\times `$ is the density function, $`H`$ the classical Hamiltonian, and $`\{,\}`$ the Poisson bracket. This equation can be seen as describing the flow of a fluid in phase space such that at any point moving along with the fluid, the density of the fluid remains constant. So besides giving the time-evolution, this equation also states a property of the time-evolution, namely that it conserves volume in phase space. In quantum mechanics we have the analogous von Neumann equation

$$\frac{d\rho }{dt}=i[\rho ,H]$$

where $`\rho :𝔏()`$ is the density operator as a function of time (note that here the derivative with respect to time is total instead of partial). This equation merely gives the time-evolution $`\rho (t)=\tau _t(\rho (0))`$ of the density operator, where $`\tau `$ is the time-evolution on the observable algebra here viewed as acting on the state instead of the observables. Von Neumann’s equation by itself should therefore not be regarded as a quantum mechanical analogue of Liouville’s theorem.

## 4 Poincaré recurrence in C\*-algebras

In Section 3 we proposed a quantum analogue of Liouville’s theorem for bounded quantum systems. So, by analogy with classical mechanics, these are the type of systems for which we could expect recurrence. In this section, however, we will be able to study Poincaré recurrence in the more general setting of abstract C\*-algebras.

As we shall see, the theory is surprisingly close to the usual measure theoretic setting. It therefore seems appropriate to briefly review Poincaré’s recurrence theorem and its proof. Let $`=\{1,2,3,\mathrm{}\}`$ be the positive integers. Consider a measure space $`(X,\mathrm{\Sigma },\mu )`$ with $`\mu (X)<\mathrm{}`$, and let $`T:XX`$ be a mapping such that $`\mu (T^1`$$`(S))=\mu (S)`$ for all $`S`$ in $`\mathrm{\Sigma }`$. This is merely an abstraction of Liouville’s theorem. For some $`S\mathrm{\Sigma }`$, suppose that $`\mu (ST^n(S))=0`$ for all $`n`$. For all $`n,k`$ we then have $`\mu (T^k(S)T^{(n+k)}(S))=\mu (T^k(ST^n(S)))=\mu (ST^n(S))=0`$. So $`\mu (T^m(S)T^n(S))=0`$ for all $`m,n`$ with $`mn`$. It follows that

$$\mu (X)\mu \left(\underset{k=1}{\overset{n}{}}T^k(S)\right)=\underset{k=1}{\overset{n}{}}\mu (T^k(S))=\underset{k=1}{\overset{n}{}}\mu (S)=n\mu (S).$$

Letting $`n\mathrm{}`$ it follows that $`\mu (S)=0`$. This is one form of Poincaré’s recurrence theorem, namely if $`\mu (S)>0`$, then there exists a positive integer $`n`$ such that $`\mu (ST^n(S))>0`$. It tells us that $`S`$ contains a set $`ST^n(S)`$ of positive measure which is mapped back into $`S`$ by $`T^n`$.

Note that the mapping $`g\tau (g)=gT`$ is a $``$-homomorphism of the C\*-algebra $`B_{\mathrm{}}(\mathrm{\Sigma })`$ into itself such that $`\phi (\tau (g))=\phi (g)`$ and $`\mu (ST^n(S))=\phi \left(\chi _S\tau ^n(\chi _S)\right)`$ for $`S\mathrm{\Sigma }`$, where $`\phi (g)=g𝑑\mu `$ for all $`gB_{\mathrm{}}(\mathrm{\Sigma })`$. Using this notation Poincaré’s recurrence theorem can be stated as follows: If $`\phi (\chi _S)>0`$, then there exists a positive integer $`n`$ such that $`\phi \left(\chi _S\tau ^n(\chi _S)\right)>0`$. The general C\*-algebraic approach will now be modelled after this situation. We also get some inspiration from Postulate 2.3, for reasons to be explained in Section 5.

Definition 4.1. Let $`𝔄`$ be a $``$-algebra, and $`𝔅`$ a unital C\*-algebra. Let $`\phi :𝔄𝔅`$ be a positive mapping (i.e. $`\phi (A^{}A)0`$ for all $`A𝔄`$). We call $`\phi `$ additive if

$$\underset{k=1}{\overset{n}{}}\phi \left(P_k\right)1$$

for any projections $`P_1,\mathrm{},P_n𝔄`$ for which $`\phi (P_kP_lP_k)=0`$ if $`k<l`$. We call $`\phi `$ faithful if it is linear, $`𝔄`$ is unital, $`\phi (1)=1`$, and $`\phi (A^{}A)>0`$ for all non-zero $`A`$ in $`𝔄`$. We call $`\phi `$ a C\*-trace if it is linear, $`𝔄`$ is unital, $`\phi (1)=1`$, and for all $`A,B𝔄`$ we have $`\phi (AB)=\phi (BA)`$. (Remember: Any C\*-algebra is a $``$-algebra.)

If the positive mapping $`\phi `$ given in Definition 4.1 is faithful, then it is also additive, as we now show. Let $`P_1,\mathrm{},P_n𝔄`$ be any projections for which $`\phi (P_kP_lP_k)=0`$ if $`k<l`$. For $`k<l`$ we then have $`\phi \left((P_lP_k)^{}P_lP_k\right)=0`$, so $`P_lP_k=0`$, and therefore $`P_kP_l=(P_lP_k)^{}=0`$. This implies that

$$\underset{k=1}{\overset{n}{}}P_k1$$

since the left-hand side is a projection in $`𝔄`$. Thus

$$\underset{k=1}{\overset{n}{}}\phi \left(P_k\right)=\phi \left(\underset{k=1}{\overset{n}{}}P_k\right)\phi (1)=1$$

as promised.

In the measure theoretic setting described above, we can assume without loss of generality that $`\mu (X)=1`$. Then $`\phi :B_{\mathrm{}}(\mathrm{\Sigma })`$ is an additive C\*-trace since

$$\underset{k=1}{\overset{n}{}}\phi \left(\chi _{S_k}\right)=\underset{k=1}{\overset{n}{}}\mu (S_k)=\mu \left(\underset{k=1}{\overset{n}{}}S_k\right)\mu (X)$$

for any $`S_1,\mathrm{},S_n\mathrm{\Sigma }`$ such that $`\phi \left(\chi _{S_k}\chi _{S_l}\right)=\mu \left(S_kS_l\right)=0`$ if $`kl`$.

We now state and prove a C\*-algebraic version of Poincaré’s recurrence theorem:

Theorem 4.2. Consider a $``$-algebra $`𝔄`$ and a unital C\*-algebra $`𝔅`$, and let $`\phi :𝔄𝔅`$ be an additive mapping. Let $`\tau :𝔄𝔄`$ be a $``$-homomorphism such that $`\phi (\tau (PQP))=\phi (PQP)`$ for all projections $`P,Q𝔄`$. Then, for any projection $`P𝔄`$ such that $`\phi (P)>0`$, there exists a positive integer $`n`$ such that $`\phi (P\tau ^n(P)P)>0`$.

Proof. Note that $`\phi (P\tau ^n(P)P)=\phi \left((\tau ^n(P)P)^{}\tau ^n(P)P\right)0`$ for all $`n`$. We now imitate the measure theoretic proof.

Suppose $`\phi (P\tau ^n(P)P)=0`$ for all $`n`$. For all $`k,n`$ we then have

$$\phi \left(\tau ^k(P)\tau ^{n+k}(P)\tau ^k(P)\right)=\phi \left(\tau ^k\left(P\tau ^n(P)P\right)\right)=\phi \left(P\tau ^n(P)P\right)=0.$$

Since $`\phi `$ is additive, it follows for any $`n`$ that

$$\underset{k=1}{\overset{n}{}}\phi \left(\tau ^k(P)\right)1.$$

Furthermore,

$$\underset{k=1}{\overset{n}{}}\phi \left(\tau ^k(P)\right)=\underset{k=1}{\overset{n}{}}\phi \left(P\right)=n\phi (P)0$$

since $`\phi `$ is positive and $`P=P^{}P`$. Hence $`0n\phi (P)1`$, and therefore $`n\phi (P)1`$. Letting $`n\mathrm{}`$, it follows that $`\phi (P)=0`$. $`\mathrm{}`$

It is clear that the measure theoretic Poincaré recurrence theorem stated above is just a special case of Theorem 4.2, since the projections of the C\*-algebra $`B_{\mathrm{}}(\mathrm{\Sigma })`$ are exactly the characteristic functions $`\chi _S`$, where $`S\mathrm{\Sigma }`$.

Note that the trace of a finite von Neumann algebra is a faithful C\*-trace, hence we have the following corollary of Theorem 4.2, which will be used in Section 5:

Corollary 4.3. Consider a finite von Neumann algebra $`𝔐`$, and let tr be its trace. Let $`\tau :𝔐𝔐`$ be a $``$-homomorphism such that tr$`(\tau (A))=`$ tr$`(A)`$ for all $`A`$ in $`𝔐`$. Then, for any projection $`P𝔐`$ such that tr$`(P)>0`$, there exists a positive integer $`n`$ such that tr$`(P\tau ^n(P))>0`$.

We can also give a C\*-algebraic version of Khintchine’s theorem (see , for example, as well as ). But first we mention that a subset $`E`$ of$``$ is called relatively dense in $``$ if there is an $`n`$ such that the set

$$E\{j,j+1,\mathrm{},j+n1\}$$

is non-void for every $`j`$.

Theorem 4.4. Consider a unital C\*-algebra $`𝔄`$, and let $`\phi :𝔄`$ be a C\*-trace. Let $`\tau :𝔄𝔄`$ be a $``$-homomorphism such that $`\tau (1)=1`$ and $`\phi (\tau (A^{}A))\phi (A^{}A)`$ for every $`A`$ in $`𝔄`$. For any projection $`P`$ in $`𝔄`$,  and any $`\epsilon >0`$, it then follows that the set

$$E=\{k:\phi (P\tau ^k(P))>\phi (P)^2\epsilon \}$$

is relatively dense in $``$.

Proof. Let $`(,\pi ,\mathrm{\Omega })`$ be the cyclic representation of $`𝔄`$ obtained from $`\phi `$ by the GNS-construction (just as from $`\omega `$ in (1)). This gives us a linear function

$$\iota :𝔄:A\pi (A)\mathrm{\Omega }$$

such that $`\iota (𝔄)`$ is dense in $``$, and

$$\iota (A),\iota (B)=\mathrm{\Omega },\pi (A^{}B)\mathrm{\Omega }=\phi (A^{}B)$$

for all $`A,B𝔄`$. For any $`A𝔄`$ we therefore have

$`\iota (\tau (A))^2`$ $`=\phi \left(\left(\tau (A)\right)^{}\tau (A)\right)=\phi \left(\tau (A^{}A)\right)`$
$`\phi (A^{}A)`$
$`=\iota (A)^2.`$
By the linearity of $`\iota `$ and $`\tau `$ it now follows that

$$\overline{\tau }:\iota (𝔄):\iota (A)\iota (\tau (A))$$

is well-defined (namely if $`\iota (A)=\iota (B)`$, then $`\iota (\tau (A))=\iota (\tau (B))`$), linear and bounded, with $`\overline{\tau }1`$. Since $`\overline{\tau }`$ is bounded, we can extend it linearly to the whole of $``$, keeping $`\overline{\tau }1`$. We are now in a position to imitate the proof of the measure theoretic Khintchine theorem.

Let $`Q`$ be the projection of $``$ onto $`\{x:\overline{\tau }x=x\}`$. For $`k`$ we have

$$\phi (P\tau ^k(P))=\iota (P),\iota (\tau ^k(P))=\iota (P),\overline{\tau }^k\iota (P)=x,\overline{\tau }^kx$$

where $`x=\iota (P)`$. By the mean ergodic theorem we know that there exists an $`n`$ such that

$$\frac{1}{n}\underset{k=0}{\overset{n1}{}}\overline{\tau }^kxQx\frac{\epsilon }{x+1}.$$

Since $`\overline{\tau }Qx=Qx`$, it follows for any $`j`$ that

$`{\displaystyle \frac{1}{n}}{\displaystyle \underset{k=j}{\overset{n+j1}{}}}\overline{\tau }^kxQx`$ $`=\overline{\tau }^j\left({\displaystyle \frac{1}{n}}{\displaystyle \underset{k=0}{\overset{n1}{}}}\overline{\tau }^kxQx\right){\displaystyle \frac{1}{n}}{\displaystyle \underset{k=0}{\overset{n1}{}}}\overline{\tau }^kxQx`$
$`{\displaystyle \frac{\epsilon }{x+1}}\text{,}`$
so

$$\left|x,\frac{1}{n}\underset{k=j}{\overset{n+j1}{}}\overline{\tau }^kxQx\right|x\frac{1}{n}\underset{k=j}{\overset{n+j1}{}}\overline{\tau }^kxQx<\epsilon .$$

Also, $`x,\iota (1)=x,Q\iota (1)=Qx,\iota (1)`$ since $`\overline{\tau }\iota (1)=\iota (\tau (1))=\iota (1)`$, so

$$\phi (P)^2=\left|\iota (P),\iota (1)\right|^2=\left|x,\iota (1)\right|^2Qx^2\iota (1)^2=x,Qx$$

since $`\iota (1),\iota (1)=\phi (1^{}1)=1`$. Thus

$$\left|\frac{1}{n}\underset{k=j}{\overset{n+j1}{}}\phi (P\tau ^k(P))\right|=\left|\frac{1}{n}\underset{k=j}{\overset{n+j1}{}}x,\overline{\tau }^kx\right|>\left|x,Qx\right|\epsilon \phi (P)^2\epsilon .$$

Since $`\phi `$ is a C\*-trace, we have $`\phi (P\tau ^k(P))=\phi ((P\tau ^k(P))^{}P\tau ^k(P))0`$ for all $`k`$, hence

$$\underset{k=j}{\overset{n+j1}{}}\phi (P\tau ^k(P))>n\left(\phi (P)^2\epsilon \right).$$

This implies that $`\phi (P\tau ^k(P))>\phi (P)^2\epsilon `$ for some $`k\{j,j+1,\mathrm{},n+j1\}`$, i.e. $`E`$ is relatively dense in $``$. $`\mathrm{}`$

Recall that a finite factor is a finite von Neumann algebra $`𝔐`$ which is also a factor, i.e. $`𝔐𝔐^{}=1`$. In this case we can therefore take the trace of $`𝔐`$ to be complex valued, so the conditions of Theorem 4.4 are satisfied when $`𝔄`$ is a finite factor and $`\phi `$ is its trace.

We also mention that if in Theorem 4.4 we consider the special case where $`𝔄`$, $`\tau `$ and $`\phi `$ are taken as $`B_{\mathrm{}}(\mathrm{\Sigma })`$, $`\tau `$ and $`\phi `$ as defined above Definition 4.1, with $`\mu (X)=1`$, then we get the usual measure theoretic theorem of Khintchine, namely, given any $`\epsilon >0`$, the set $`\{k:\mu (ST^k(S))>\mu (S)^2\epsilon \}`$ is relatively dense in $``$ for all $`S\mathrm{\Sigma }`$, where the condition $`\mu (T^1(S))=\mu (S)`$ can now be weakened to $`\mu (T^1(S))\mu (S)`$. This is a stronger result than Poincaré’s recurrence theorem (in the form stated above), despite the slightly weaker assumptions.

## 5 Physical interpretation

Consider a bounded quantum system $`(𝔐,,H)`$ and assume that $`𝔐`$ is a factor. Let $`\tau `$ be the system’s time-evolution, as in Proposition 3.3. Fix any $`t>0`$. Since the trace tr of $`𝔐`$ is faithful, Corollary 4.3 and Proposition 3.3 tell us that for any non-zero projection $`P𝔐`$ there exists an $`n(t)`$ such that

$$\text{tr}\left(P\tau _{n(t)t}(P)\right)>0.$$
(5)
Note that tr$`(P\tau _{n(t)t}(P))=`$ tr$`(P\tau _{n(t)t}(P)P)`$, which is similar to the form of $`\omega ^{}`$ in Postulate 2.3, i.e. a state after a “yes” was obtained in a yes/no experiment with projection $`P`$. We now look at this similarity more closely by exploiting the analogy between quantum and classical mechanics described in Sections 2 and 3.

In Section 3 we saw that tr can be viewed as a quantum analogue of integration over a bounded set in phase space with respect to Lebesgue measure $`\lambda `$. In order to apply Poincaré’s recurrence theorem to classical mechanics, we know that the system in question has to be confined to a bounded (Borel) set $`F`$ in the phase space $`^{2n}`$, i.e. $`\lambda (F)<\mathrm{}`$. Since we can assume without loss that $`\lambda (F)>0`$, we can normalize $`\lambda `$ on $`F`$ by defining a measure $`\lambda ^{}`$ on the Borel sets of $`^{2n}`$ by

$$\lambda ^{}(S)=\lambda (SF)/\lambda (F)\text{.}$$

If we now view $`\lambda ^{}`$ as describing a state of the system (as explained in Section 2), then it essentially says that every part of $`F`$ is equally likely to contain the state of the system (viewed as a point in the phase space). In other words, when we know nothing about the state of the system (aside from the fact that it is in $`F`$), then we can describe it by $`\lambda ^{}`$, or in C\*-algebraic terms by the state $`\phi `$ on $`B_{\mathrm{}}(^{2n})`$ defined by

$$\phi (g)=g𝑑\lambda ^{}.$$

Since tr$`(1)=1`$ and $`𝔐`$ is a factor, tr is a state on $`𝔐`$, and therefore we view tr as the quantum analogue of $`\phi `$. By this analogy we would expect tr to describe the state of our quantum system when we know nothing about the system’s state. This is indeed true in the special case where $``$ is finite dimensional and $`𝔐=𝔏()`$, since for any rank one projection $`Q`$ in $`𝔐`$ we then have tr$`(Q)=1/dim()`$ which tells us that all values are equally probable when we measure an observable (assuming the observable has no degenerate eigenvalues). Furthermore, since tr is ultraweakly continuous, it is a normal state and hence it is given by a density operator (see ), as one would expect for a physically meaningful state. We therefore suggest the following hypothesis:

Postulate 5.1. Consider a bounded quantum system $`(𝔐,,H)`$ where $`𝔐`$ is a factor. When we have no information regarding the state of the system, the state is given by the trace tr of $`𝔐`$.

So look at the case where we have no information about the state of our bounded quantum system. By Postulate 5.1 the state is then given by tr. At time $`0`$ we perform a yes/no experiment with projection $`P𝔐`$ on the system. Assuming the result is “yes”, the state of the system after the experiment is given by the state $`\omega `$ on $`𝔐`$ defined by

$$\omega (A)=\text{ tr}(PA)/\text{tr}(P)\text{,}$$

according to Postulate 2.3. (Also recall from Section 2 that the probability of getting “yes” is tr$`(P)`$, therefore tr$`(P)>0`$ in this case.) By (5) we then have

$$p(t):=\omega (\tau _{n(t)t}(P))>0.$$
(6)
This simply tells us that if we were to repeat the above mentioned yes/no experiment exactly at the moment $`n(t)t`$, then there is a non-zero probability $`p(t)`$ that we will again get “yes”. By replacing $`t`$ by $`t^{}=n(t)t+1`$, we see that there is in fact an unbounded set of moments $`n(t)t<n(t^{})t^{}<\mathrm{}`$ for which (6) holds.

So we have obtained a quantum mechanical version of recurrence. Note that the measure theoretic Poincaré recurrence theorem stated in Section 4 will give exactly the same result as (6), with the same physical interpretation, when applied to classical mechanics; just replace $`\omega `$, tr, $`\tau `$ and $`P`$ by their classical analogues described in Section 2 and in this section. So we see that recurrence in quantum mechanics and in classical mechanics follow from the same theorem, namely Theorem 4.2, since Corollary 4.3 and measure theoretic Poincaré recurrence are both special cases of this theorem.

Of course, Theorem 4.4 tells us that for any $`\epsilon >0`$ there is in fact a relatively dense set $`M`$ in $``$ such that

$$\omega (\tau _{mt}(P))>\text{ tr}(P)\epsilon $$
(7)
for all $`mM`$. Since tr$`(P)`$ was the probability of getting a “yes” during the first execution of the yes/no experiment, we see from (7) that at the moments $`mt`$ the probability of getting “yes” when doing the experiment a second time is larger or at least arbitrarily close to the original probability of getting “yes”. Similar results concerning wave functions and density operators are presented in  and . If as before we replace $`\omega `$, tr, $`\tau `$ and $`P`$ by their classical counterparts, and then apply Theorem 4.4 again, we find the same result as (7) for classical mechanics, with exactly the same interpretation as in quantum mechanics.

There is, however, a small technical problem: The probability of repeating the yes/no experiment exactly at the moment $`n(t)t`$ is zero. The same goes for any of the moments $`mt`$ above. The next simple proposition remedies the situation in the quantum case:

Proposition 5.2. Let $`\tau `$ be as in Proposition 3.3, where we take $`𝔐`$ to be a finite factor. Then for any projection $`P`$ in $`𝔐`$, the mapping

$$:t\text{ tr}(P\tau _t(P))$$

is continuous, where tr is the trace of $`𝔐`$.

Proof. By Stone’s theorem $`U_t`$ in Proposition 3.3 is strongly continuous, so clearly the mapping $`t\tau _t(A)`$ is weakly continuous for every $`A\dot{𝔐}`$. Hence $`tP\tau _t(P)`$ is weakly continuous. We know that tr is ultraweakly continuous (see , for example) and therefore it is weakly continuous on the unit ball. Since $`P\tau _t(P)1`$, we conclude that $`t`$ tr$`(P\tau _t(P))`$ is continuous. $`\mathrm{}`$

So from (7) we see that for every $`mM`$ there exists a $`\delta _m>0`$ such that

$$\omega (\tau _s(P))>\text{tr}(P)\epsilon \text{for }mt\delta _m<s<mt+\delta _m.$$

This tells us that quantum mechanical recurrence is possible in practice, assuming we are working with a bounded quantum system as above, since there is a non-zero probability of repeating the yes/no experiment during one of the time-intervals $`(mt\delta _m,mt+\delta _m)`$.

Of course, this remark leads to the next question: Which physical systems can be mathematically described as bounded quantum systems with the observable algebras being factors?

In classical mechanics Poincaré’s recurrence theorem applies to systems that are confined to a bounded set in phase space. From a physical standpoint this is true if the system is confined to a finite volume in space, and it is isolated from outside influences (which could increase its energy content), to prevent any of its momentum components to go to infinity. (To see this, use Cartesian coordinates. Here we assume that each potential of the form $`1/r`$ or the like has some “cut-off” at small values of $`r`$, since for example particles are of finite size and collide when they get too close, the point being that there is not an infinite amount of potential energy available in the system.)

Most likely then (keeping in mind the close analogy between quantum and classical mechanics), recurrence will occur for quantum mechanical systems bounded in space and isolated from outside influences (apart from the yes/no experiments we perform on it). This is confirmed by  and . So we might guess that these types of systems can be described as bounded quantum systems in the sense of Definition 3.2 with $`𝔐`$ a factor. This seems to be related to the nuclearity requirement in quantum field theory (see ), where a bounded set in classical phase space is intuitively thought of as corresponding to a finite dimensional subspace of the quantum state space. Since a quantum system whose state space $``$ is finite dimensional is clearly a bounded quantum system (since $`𝔏()`$ is a finite factor in this case), it certainly does not seem too far-fetched to conjecture that a physical system bounded in space and isolated from outside influences can be mathematically described as a bounded quantum system with a factor as the observable algebra. We will not pursue these matters further in this paper however.

Acknowledgements

I would like to thank Prof. A. Ströh for encouraging me to pursue this line of research, and the National Research Foundation for financial support.

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