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# Linear stochastic differential equations with functional boundary conditions

## 1 Introduction

It is well known that, under suitable Lipschitz and growth conditions on the coefficients, a classical Itô stochastic differential equation

(1.1)
$$X(t)=\xi +_0^tb(s,X(s))𝑑s+_0^t\sigma (s,X(s))𝑑W(s),$$

where $`W`$ is a Wiener process and $`\xi `$ is a $`_0`$-measurable random variable for a given non-anticipating filtration $`\{_t,t0\}`$ of $`W`$, has a unique strong solution which is a Markov process.

If $`\xi `$ is not $`_0`$-measurable or the coefficients $`b`$, $`\sigma `$ are random and non-adapted, then any reasonable interpretation of $`X`$ in (1.1) will not be an $`_t`$-adapted process and, unless $`\sigma `$ is a constant, we need to use some anticipating stochastic integral to give a sense to the equation. In these cases, the solution is not a Markov process in general.

Still another setting that leads to anticipation is the case of boundary conditions. That means, the first variable of the solution process is no longer a datum of the problem, time runs in a bounded interval, say from 0 to 1, and we impose a relation $`h(X(0),X(1))=0`$ between the first and the last variables of the solution. In this situation, the fact that the solution will not be Markovian is quite intuitive, since the strong relationship between $`X(0)`$ and $`X(1)`$ will prevent the independence of $`X(0)`$ and $`X(1)`$ from holding, even when conditioning to $`X(a)`$, $`a]0,1[`$, except maybe in some very particular cases.

On the other hand, it may also seem intuitive that the following weaker conditional independence property can hold true: For any $`0a<b1`$, the $`\sigma `$-fields $`\sigma \{X(t),t[a,b]\}`$ and $`\sigma \{X(t),t]a,b[^c\}`$ are conditionally independent given $`\sigma \{X(a),X(b)\}`$. We will denote it by

(1.2)
$$\sigma \{X(t),t[a,b]\}\underset{\sigma \{X(a),X(b)\}}{\text{ }\text{ }\text{ }\text{ }\text{ }}\sigma \{X(t),t]a,b[^c\}.$$

Now $`X(0)`$ and $`X(1)`$ are on the same side in relation (1.2), so that the boundary condition does not seem to cause the problem seen above. But the following example shows that this is wrong:

###### Example 1.1

Consider the problem

$$\{\begin{array}{c}\dot{X}(t)=f(X(t))dt+\dot{W}(t),t[0,1]\hfill \\ h(X(0),X(1))=0,\hfill \end{array}$$

where the noise appears additively, and assume that a unique solution exists and that the boundary condition given by $`h`$ does not reduce to an initial or final condition. Then, relation (1.2) holds if and only if $`f(x)=\alpha x+\beta `$, for some constants $`\alpha `$ and $`\beta `$. This was proved in Nualart and Pardoux .

The processes satisfying (1.2) were called reciprocal processes by S. Bernstein . The concept arose directly from E. Schrödinger ideas on the formulation of quantum mechanics. More recent research on such processes has been carried out by B. Jamison , A. Krener , R. Frezza, A. Krener and B. Levy , M. Thieullen  and J.C. Zambrini .

Other names can be found in the literature to refer to the same concept. A reciprocal process is a one-parameter Markov field in Paul Lévy’s terminology, and is also called a quasi-Markov process, a local Markov process and a Bernstein process. We shall simply call them Markov fields (see Definition 4.1).

###### Example 1.2

Consider now the problem

$$\{\begin{array}{c}\ddot{X}(t)+f(X(t),\dot{X}(t))=\dot{W}(t),t[0,1]\hfill \\ X(0)=c_1,X(1)=c_2.\hfill \end{array}$$

This is a second order stochastic differential equation, and it is natural to ask for conditional independence properties of the 2-dimensional process $`Y(t)=(\dot{X}(t),X(t))`$, since $`X(t)`$ has $`C^1`$ paths, and therefore it is meaningless to look for this kind of properties for $`X(t)`$ itself.

Nualart and Pardoux  proved that if $`Y(t)`$ is a Markov field, then, as in Example 1.1, $`f`$ must be an affine function. Moreover, if $`f`$ is affine, then $`Y`$ is not only a Markov field, but a Markov process.

Let us look at this example more closely: Note that a boundary condition for a second order equation has the general form

$$h(Y(0),Y(1))=(0,0),$$

where $`Y(0)=(\dot{X}(0),X(0))`$ and $`Y(1)=(\dot{X}(1),X(1))`$. However, in Example 1.2 the two scalar conditions do not mix values at 0 and values at 1 of $`Y`$. The same happens, for instance, with the Neumann-type conditions $`\dot{X}(0)=c_1`$, $`\dot{X}(1)=c_2`$, and the result is the same ($`Y`$ Markov field $``$ $`f`$ affine $``$ $`Y`$ Markov process).

From these examples and other equations of first and second order that have been studied so far (see e.g. , , , ), we learn that

1. The Markovian properties can be expected only in “linear” cases.
2. The specific Markovian property depends on the actual form of the boundary condition.

It should also be noted that the requirement of linearity on the drift coefficient $`f`$ is related to the fact that the noise appears additively. Should not this be the case, the Markovian property would occur under a different condition which relates the drift and the diffusion coefficients (see  and ).

In the present paper we will consider linear stochastic differential equations of arbitrary order with additive white noise. Our boundary conditions will not be restricted to involve the solution process at the endpoints of the time interval, but we will allow them to involve the values at finitely many points inside the interval. They are usually called functional or lateral boundary conditions. Our main goal is to seek which kind of conditional independence properties can be established for the solution. A preliminary work in this direction was published in Alabert and Ferrante . Here we considerably refine and extend the results therein. This type of equations was already considered by Russek , who proved that the solutions are Markov processes if and only if the lateral conditions fix to a constant the variables $`X(t)`$, for all points $`t`$ in the support of the conditions. His techniques, based in the notion of reproducing kernel space, are different from ours.

Our main result (Theorem 4.5) can be stated in the following way: Fix two points $`0a<b1`$ and set $`Y(t)=(D^{n1}X(t),\mathrm{},DX(t),X(t))`$, where $`n`$ is the order of the equation, $`X(t)`$ is its solution process, and $`D`$ is the time derivative; the process $`\{Y(t),t[0,1]\}`$ satisfies the relation (1.2) if and only if there are no lateral conditions involving points inside and outside the interval $`[a,b]`$. We also state a conditional independence property for the case when there are conditions that do involve points inside and outside $`[a,b]`$ (Theorem 4.10). Finally we obtain a result from which Russek’s theorem can be trivially recovered (Theorem 4.11). The paper is organised as follows:

In Section 2 we precise the statement of the problem and develop some notation and properties that will be needed later.

Section 3 contains the main probabilistic tools: Lemma 3.4 and Proposition 3.7. The first is a characterisation of the conditional independence of two random vectors given a function of them. It is the most important ingredient in the proof of Theorem 4.5, but it cannot be applied for certain singular values of $`a`$ and $`b`$. For these values, we employ an approximation by the solution of perturbed equations. The approximation argument involves the convergence in $`L^2`$ of a sequence of conditional expectations with varying conditioning $`\sigma `$-fields. Proposition 3.7 gives a sufficient condition for this convergence in a general setting.

In Section 4 we establish the main results. The proofs of Theorems 4.10 and 4.11 will be only sketched, since the procedure is similar to that of Theorem 4.5, with slight modifications.

## 2 Linear SDE with functional boundary conditions

The present section will be devoted to the statement of the problem, the definition of a solution, and to absolute continuity and approximation results for the solution of an $`n`$–th order linear stochastic differential equation with linear functional boundary conditions.

### 2.1 Statement of the problem and definition of a solution

Consider the differential operator

$$L:=D^n+a_{n1}D^{n1}+\mathrm{}+a_1D+a_0,D:=\frac{d}{dt},$$

where $`a_i`$ are continuous functions on $`[0,1]`$. Let $`\{W(t),t[0,1]\}`$ be a standard Wiener process. We assume that $`W`$ is the coordinate process in the classical Wiener space $`(\mathrm{\Omega },,P)`$, that means, $`\mathrm{\Omega }=C_0([0,1];IR)`$ is the space of continuous functions on $`[0,1]`$ vanishing at zero, $``$ its Borel $`\sigma `$-field, and $`P`$ the Wiener measure. We shall deal with the SDE

(2.1)
$$L[X]=\dot{W}$$

on $`[0,1]`$, together with the additional conditions

(2.2)
$$\underset{j=1}{\overset{m}{}}\alpha _{ij}X(t_j)=c_i,1in,$$

where $`mn`$, $`0t_1<\mathrm{}<t_m1`$ are some given points in $`[0,1]`$, and $`\alpha _{ij}`$, $`c_i`$ are real numbers. The matrix of coefficients $`(\alpha _{ij})`$ is assumed to have full rank.

As in the case of ordinary differential equations, (2.1)-(2.2) can be regarded as a first order system

(2.3)
$$DY(t)+A(t)Y(t)=\dot{B}(t),t[0,1],$$

with constraints

(2.4)
$$\underset{j=1}{\overset{m}{}}\alpha _{ij}Y_n(t_j)=c_i,1in,$$

where $`Y(t)=(Y_1(t),\mathrm{},Y_n(t))`$, $`Y_i(t)=D^{ni}X(t)`$ for $`1in`$, $`B(t)=(W(t),0,\mathrm{},0)`$, and

(2.5)
$$A(t)=\left[\begin{array}{cccccc}a_{n1}(t)& a_{n2}(t)& \mathrm{}& a_1(t)& a_0(t)& \\ 1& 0& \mathrm{}& 0& 0& \\ 0& 1& \mathrm{}& 0& 0& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ 0& 0& \mathrm{}& 1& 0& \end{array}\right].$$

The lateral condition (2.4) is a special case of the general linear condition

(2.6)
$$\mathrm{\Lambda }[Y]=c,$$

for $`\mathrm{\Lambda }`$ in the set $`(C([0,1];IR^n);IR^n)`$ of linear continuous $`IR^n`$-valued functionals on $`C([0,1];IR^n)`$, and $`cIR^n`$. By the Riesz representation theorem, (2.6) can be written as

(2.7)
$$_0^1𝑑F(t)Y(t)=c,$$

where $`F`$ is an $`(n\times n)`$-matrix whose components are functions of bounded variation.

When the right-hand side of (2.3) is a continuous vector function $`g`$, it is well known that the system

(2.8)
$$\{\begin{array}{c}DY(t)+A(t)Y(t)=g(t),t[0,1]\hfill \\ _0^1𝑑F(t)Y(t)=c\hfill \end{array}$$

admits a unique solution, which belongs to $`C^1([0,1];IR^n)`$, if and only if for some $`s[0,1]`$ (equivalently, for every $`s[0,1]`$)

$$\text{(H0)}det_0^1𝑑F(t)\mathrm{\Phi }^s(t)0,$$

where $`\mathrm{\Phi }^s(t)`$ denotes the fundamental matrix solution of $`DY(t)+A(t)Y(t)=0`$, that is, $`s[0,1]`$,

$$\{\begin{array}{c}\frac{d}{dt}\mathrm{\Phi }^s(t)+A(t)\mathrm{\Phi }^s(t)=0,t[0,1]\hfill \\ \mathrm{\Phi }^s(s)=\text{I},\hfill \end{array}$$

with I the identity matrix. In turn, this is equivalent to say that the homogeneous problem ($`g0`$, $`c0`$) has only the trivial solution. When hypothesis (H0) holds, the solution to (2.8) is given by

$$Y(t)=J(t)^1c+_0^1G(t,s)g(s)𝑑s,$$

where

(2.9)
$$J(t)=_0^1𝑑F(u)\mathrm{\Phi }^t(u)$$

and $`G(t,s)`$ is the (matrix-valued) Green function associated to $`A`$ and $`F`$. An explicit expression for this function is the following (see e.g.  or ):

(2.10)
$$G(t,s)=J(t)^1\left[_0^s𝑑F(u)J(u)^1\text{1}_{\{ts\}}\text{I}\right]J(s).$$

Under (H0), we define the solution to (2.3)-(2.4) as the $`n`$-dimensional stochastic process

(2.11)
$$Y(t)=J(t)^1c+_0^1G(t,s)𝑑B(s),$$

and the solution to (2.1)-(2.2) as the process $`\{X(t)=Y_n(t),t[0,1]\}`$. The Green function (2.10) has bounded variation, so that the Wiener integrals in (2.11) can be interpreted pathwise by means of an integration by parts

$$\left[_0^1G(t,s)𝑑B(s)\right](\omega )=_0^1G(t,ds)B(s)(\omega )$$

(we take into account here that $`G(t,1)=0`$, $`t`$), and therefore $`Y`$ can be defined everywhere. We shall assume throughout the paper that the solution is interpreted in this pathwise sense. Furthermore, it is not difficult to verify that the process $`Y(t)`$ so defined is continuous (hence $`X(t)`$ is a $`C^{n1}`$ process) and that, for each $`t[0,1]`$, the mapping $`\omega Y(\omega )`$ from $`\mathrm{\Omega }`$ into $`C([0,1];IR^n)`$ is continuous with the usual topologies.

Notice that, with the notation introduced in (2.7), the particular lateral condition (2.4) corresponds to

(2.12)
$$dF=\left[\begin{array}{ccccc}0& \mathrm{}& 0& \underset{j=1}{\overset{m}{}}\alpha _{1j}\delta _{t_j}& \\ 0& \mathrm{}& 0& \underset{j=1}{\overset{m}{}}\alpha _{2j}\delta _{t_j}& \\ \mathrm{}& & \mathrm{}& \mathrm{}& \\ 0& \mathrm{}& 0& \underset{j=1}{\overset{m}{}}\alpha _{nj}\delta _{t_j}& \end{array}\right],$$

where $`\delta _t`$ denotes the Dirac measure at $`t`$, and that $`J_{ik}(t)=_{j=1}^m\alpha _{ij}\mathrm{\Phi }_{nk}^t(t_j)`$. Notice also that only the first column of $`G(t,s)`$ is relevant in (2.11).

Another natural definition of solution for the system (2.3)-(2.4) arises if, for each $`\omega `$ fixed, we consider the object $`\dot{B}(\omega )`$ as the derivative of a continuous $`IR^n`$-valued function defined on $`[0,1]`$, and therefore we regard (2.3) as an equation between distributions. The vector function $`Y=(Y_1,\mathrm{},Y_n)`$ will be a solution in the distributional sense if for any smooth vector $`\phi =(\phi _1,\mathrm{},\phi _n)`$ vanishing in the complement of $`]0,1[`$, with $`_0^1\phi =0`$,

(2.13)
$$_0^1\left(Y(t)+_0^tA(a)Y(s)𝑑sB(t)\right)\phi (t)𝑑t=0$$

and (2.4) is satisfied. But (2.13) amounts to say that there exists a constant $`Y(0)`$ such that

$$Y(t)Y(0)+_0^tA(s)Y(s)𝑑s=B(t),t[0,1],$$

and a fortiori we find that $`Y`$ must be a continuous function. It is easily seen that both concepts of solution coincide.

### 2.2 On the law of the solution

In the present subsection we shall prove an absolute continuity result for the law of the solution process $`\{Y(t),t[0,1]\}`$. Here we allow the boundary condition (2.4) to depend on all coordinates of $`Y`$, since we will use this generality later on.

Let $`Y=(Y_1,\mathrm{},Y_n):\mathrm{\Omega }C([0,1];IR^n)`$ be the solution to the problem

(2.14)
$$\{\begin{array}{c}DY(t)+A(t)Y(t)=\dot{B}(t),t[0,1]\hfill \\ \mathrm{\Lambda }[Y]=c,\hfill \end{array}$$

where $`\mathrm{\Lambda }`$ is any linear operator on $`C([0,1];IR^n)`$ with finite support $`\text{supp}\mathrm{\Lambda }=\{t_1,\mathrm{},t_m\}`$. (We are not assuming here that $`\mathrm{\Lambda }`$ involves only the coordinate function $`Y_n`$, but we do assume that problem (2.14) is well-posed.)

If $`\{s_1,\mathrm{},s_k\}[0,1]`$ is a set containing $`\text{supp}\mathrm{\Lambda }`$, then $`\mathrm{\Lambda }`$ can be regarded as a linear operator on the space of functions $`\left(\{s_1,\mathrm{},s_k\}IR^n\right)IR^{n\times k}`$. We keep the same symbol $`\mathrm{\Lambda }`$ for both interpretations. Denote by $`M`$ the linear manifold in $`IR^{n\times k}`$:

(2.15)
$$M:=\{x=(\left(\begin{array}{c}x_{11}\\ \mathrm{}\\ x_{1n}\end{array}\right),\mathrm{},\left(\begin{array}{c}x_{k1}\\ \mathrm{}\\ x_{kn}\end{array}\right))IR^{n\times k}:\mathrm{\Lambda }[x]=c\}.$$

###### Proposition 2.1

With the notations above, the random vector $`(Y(s_1),\mathrm{},Y(s_k)):\mathrm{\Omega }IR^{n\times k}`$ is absolutely continuous with respect to the Hausdorff measure in $`M`$.

Proof: Taking into account that the vector $`(Y(s_1),\mathrm{},Y(s_k))`$ is Gaussian, it suffices to prove that any open ball in $`M`$ has a positive probability under the law of this vector.

Fix $`xM`$. Let us see first that there exists $`\omega \mathrm{\Omega }`$ such that the function $`Y(\omega ):[0,1]IR^n`$ satisfies $`(Y(\omega )(s_1),\mathrm{},Y(\omega )(s_k))=x`$. Indeed, by simple interpolation, there obviously exists a $`C^{\mathrm{}}`$ function $`y:[0,1]IR^n`$ such that $`y_iy_{i+1}^{}`$, $`i=1,\mathrm{},n1`$, and $`(y(s_1),\mathrm{},y(s_k))=x`$ (therefore $`\mathrm{\Lambda }[y]=c`$). Defining

(2.16)
$$\omega (t)=y^1(t)y^1(0)+_0^t(a_{n1}(s)y^1(s)+\mathrm{}+a_0(s)y^n(s))𝑑s,$$

we find that $`y`$ is the solution path $`Y(\omega )`$ of (2.14). Any open ball $`U(x)`$ of $`M`$ centred at a point $`xM`$ has therefore a non-empty inverse image $`Y^1(B(x))\mathrm{\Omega }`$. Moreover, since the mapping $`\omega Y(\omega )`$ is continuous, $`Y^1(U(x))`$ is open. We get that $`P\{YU(x)\}>0`$.

###### Remark 2.2

Proposition 2.1 remains valid, with a similar proof, if the domain where the problem is considered consists of two disjoint intervals, say $`[0,a]`$ and $`[b,1]`$, instead of a single one (in that case $`2n`$ lateral conditions are necessary for the problem to be well-posed). The function $`\omega `$ can be defined as in (2.16) for $`t[0,a]`$; as

$$\omega (t)=y^1(t)y^1(b)+_b^t(a_{n1}(s)y^1(s)+\mathrm{}+a_0(s)y^n(s))𝑑s$$

for $`t[b,1]`$, and arbitrarily (continuous) on $`]a,b[`$.

### 2.3 An approximation result

We shall now state an easy approximation result (Proposition 2.4) that we will need partially in the proof of Proposition 4.8. Consider the space $`C^k:=C^k([0,1];IR^{n\times n})`$, with $`k`$ a fixed non-negative integer or $`\mathrm{}`$, endowed with its natural topology. Let $`𝒞`$ be the subset of $`C^k`$ comprising the matrix functions $`A:[0,1]IR^{n\times n}`$ of the form (2.5), with the topology induced by $`C^k`$. Fix a linear operator $`\mathrm{\Lambda }:C([0,1];IR^n)IR^n`$ of rank $`n`$ and with the form given by (2.12), and consider the deterministic problems:

(2.17)
$$\{\begin{array}{c}DY(t)+A(t)Y(t)=0\hfill \\ \mathrm{\Lambda }[Y]=0\hfill \end{array}$$

with $`AC^k`$. Let $`𝒟C^k`$ the class of matrix functions $`A`$ such that (2.17) has only the trivial solution. Finally set $`V:=𝒞𝒟`$.

###### Lemma 2.3

$`V`$ is open and dense in $`𝒞`$.

Proof: For $`A𝒟`$, denote by $`\mathrm{\Phi }_A^0(t)`$ the fundamental matrix solution of the system $`DY(t)+A(t)Y(t)=0`$, with $`Y(0)=\text{I}`$. Consider the composition of linear operators

IRn

ΓA

C([0,1];IRn)

Λ

IRnc--[tΦA0(t)c]--Λ[tΦA0(t)c]IRn

ΓA

C([0,1];IRn)

Λ

IRnc--[tΦA0(t)c]--Λ[tΦA0(t)c]\begin{array}[]{l}\hbox{$I\!\!R^{n}$}\vbox{\offinterlineskip\hbox to42.67912pt{\hfil\hbox{$\Gamma_{A}$}\hfil}
\kern 0.0pt\hbox to42.67912pt{\kern 1.5pt\rightarrowfill\kern 1.5pt}}\hbox{$C([0,1];I\!\!R^{n})$}\vbox{\offinterlineskip\hbox to60.64238pt{\hfil\hbox{$\Lambda$}\hfil}
\kern 0.0pt\hbox to60.64238pt{\kern 1.5pt\rightarrowfill\kern 1.5pt}}\hbox{$I\!\!R^{n}$}\\[5.69054pt]
\hskip 4.27559pt\hbox{$c$}\hbox to52.50603pt{\kern 1.5pt$\mapstochar\mathrel{\mkern-4.0mu}\mathord{-}\mkern-6.0mu\cleaders\hbox{$\mkern-2.0mu\mathord{-}\mkern-2.0mu$}\hfill\mkern-6.0mu\mathord{\rightarrow}$\kern 1.5pt}\hbox{$[t\mapsto\Phi_{A}^{0}(t)c]$}\hbox to48.23044pt{\kern 1.5pt$\mapstochar\mathrel{\mkern-4.0mu}\mathord{-}\mkern-6.0mu\cleaders\hbox{$\mkern-2.0mu\mathord{-}\mkern-2.0mu$}\hfill\mkern-6.0mu\mathord{\rightarrow}$\kern 1.5pt}\hbox{$\Lambda[t\mapsto\Phi_{A}^{0}(t)c]$}\end{array}
The mapping

Ck

φ

(IRn;IRn)A--ΛΓACk

φ

(IRn;IRn)A--ΛΓA\begin{array}[]{l}\hbox{$C^{k}$}\vbox{\offinterlineskip\hbox to42.67912pt{\hfil\hbox{$\varphi$}\hfil}
\kern 0.0pt\hbox to42.67912pt{\kern 1.5pt\rightarrowfill\kern 1.5pt}}\hbox{${\cal L}(I\!\!R^{n};I\!\!R^{n})$}\\[5.69054pt]
\hskip 1.72707pt\hbox{$A$}\hbox to53.10165pt{\kern 1.5pt$\mapstochar\mathrel{\mkern-4.0mu}\mathord{-}\mkern-6.0mu\cleaders\hbox{$\mkern-2.0mu\mathord{-}\mkern-2.0mu$}\hfill\mkern-6.0mu\mathord{\rightarrow}$\kern 1.5pt}\hbox{$\Lambda\circ\Gamma_{A}$}\end{array}
is continuous: Indeed,

$$\mathrm{\Lambda }\mathrm{\Gamma }_A\mathrm{\Lambda }\mathrm{\Gamma }_B_{(IR^n;IR^n)}\mathrm{\Lambda }_{(C,IR^n)}\mathrm{\Gamma }_A\mathrm{\Gamma }_B_{(IR^n,C)},$$

and

$$\mathrm{\Gamma }_A\mathrm{\Gamma }_B_{(IR^n,C)}=\underset{c=1}{sup}\left(\mathrm{\Phi }_A^0(t)\mathrm{\Phi }_B^0(t)\right)c_{\mathrm{}}=\underset{i,j}{\mathrm{max}}\underset{t}{sup}|\mathrm{\Phi }_A^0(t)_{i,j}\mathrm{\Phi }_B^0(t)_{i,j}|.$$

The continuity follows from the uniform continuous dependence of the solution with respect to the data in a linear Cauchy problem.

Now we use the fact that the set $`H`$ of invertible operators on $`IR^n`$ is open in $`(IR^n,IR^n)`$. We obtain that $`\phi ^1(H)`$ is open in $`C^k`$. However $`\phi ^1(H)`$ is the set of matrices $`A`$ such that $`det(\mathrm{\Lambda }\mathrm{\Gamma }_A)0`$, which coincides with $`𝒟`$ by definition. This shows that $`𝒟`$ is open in $`C^k`$.

Since $`𝒞`$ is a linear manifold in $`C^k`$, we have that $`V=𝒞𝒟`$ is open in $`𝒞`$. Note that the particular form of $`\mathrm{\Lambda }`$ does not play any role up to this point.

To prove the density, we start by checking that $`V\mathrm{}`$. If $`A𝒞`$ and $`s[0,1]`$, the corresponding fundamental matrix $`\mathrm{\Phi }_A^s()`$ has the form

$$\mathrm{\Phi }_A^s(t)=\left[\begin{array}{cccc}D^{n1}\varphi _1(t)& \mathrm{}& D^{n1}\varphi _n(t)& \\ \mathrm{}& & \mathrm{}& \\ D\varphi _1(t)& \mathrm{}& D\varphi _n(t)& \\ \varphi _1(t)& \mathrm{}& \varphi _n(t)& \end{array}\right]$$

with

$$\mathrm{\Phi }_A^s(s)=\text{I},$$

for some $`C^{n+k}`$ real functions $`\varphi _1,\mathrm{},\varphi _n`$. Conversely, any such matrix is the fundamental matrix solution $`\mathrm{\Phi }^s`$ of $`DY(t)+A(t)Y(t)=0`$ for some $`A𝒞`$. We have

$$_0^1𝑑F(t)\mathrm{\Phi }^s(t)=\left[\begin{array}{cccc}\alpha _{11}& \mathrm{}& \alpha _{1m}& \\ \mathrm{}& & \mathrm{}& \\ \alpha _{n1}& \mathrm{}& \alpha _{nm}& \end{array}\right]\left[\begin{array}{cccc}\varphi _1(t_1)& \mathrm{}& \varphi _n(t_1)& \\ \mathrm{}& & \mathrm{}& \\ \varphi _1(t_m)& \mathrm{}& \varphi _n(t_m)& \end{array}\right].$$

Take $`s\{t_1,\mathrm{},t_m\}`$. Since $`(\alpha _{ij})`$ has full rank and $`nm`$, we can obviously find numbers $`\varphi _i(t_j)`$, $`1in`$, $`1jm`$ such that this product is an invertible square matrix. Then we take $`C^{n+k}`$ functions $`\varphi _1,\mathrm{},\varphi _n`$ interpolating these numbers and so that $`\left(D^{nj}\varphi _i(s)\right)_{i,j}=\text{I}`$. The corresponding $`A`$ will therefore belong to $`V`$.

Given now $`A𝒟`$, let us fix $`A_0V`$. For $`\lambda IR`$, define

$$M_\lambda :=\left[\begin{array}{ccccc}0& 0& \mathrm{}& 0& \\ 1& 0& \mathrm{}& 0& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ 0& \mathrm{}& 1& 0& \end{array}\right]+\left[\begin{array}{ccccc}\lambda & 0& \mathrm{}& 0& \\ 0& 0& \mathrm{}& 0& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ 0& 0& \mathrm{}& 0& \end{array}\right]A_0+\left[\begin{array}{ccccc}1\lambda & 0& \mathrm{}& 0& \\ 0& 0& \mathrm{}& 0& \\ \mathrm{}& \mathrm{}& \mathrm{}& \mathrm{}& \\ 0& 0& \mathrm{}& 0& \end{array}\right]A.$$

Let us see that the function $`u:\lambda det\mathrm{\Lambda }[\mathrm{\Phi }_{M_\lambda }^0]`$ is analytical: Indeed, $`M_\lambda `$ depends analytically on $`\lambda `$, and so the fundamental solution $`\mathrm{\Phi }_{M_\lambda }^0`$ is also analytic in $`\lambda `$. Finally, analyticity is preserved by the linear functional $`\mathrm{\Lambda }`$ and the determinant. Now assume $`u0`$ in a neighbourhood of 0. This would imply $`u0`$ on the whole line. However

$$u(1)=det\mathrm{\Lambda }[\mathrm{\Phi }_{A_0}^0]0.$$

We conclude that there exists a sequence $`\{\lambda _n\}_n`$ converging to zero such that $`M_{\lambda _n}V`$. Since $`M_{\lambda _n}A`$ as $`n\mathrm{}`$, the density is proved.

This proof borrows some ideas from Theorem 7.1 in Chow and Lasota .

###### Proposition 2.4

Let $`A^N(t)`$ be a sequence of functions in $`V`$ converging to $`A(t)V`$. Let $`Y^N(t)`$ and $`Y(t)`$ be the corresponding unique solutions to:

$$\{\begin{array}{c}DY^N(t)+A^N(t)Y^N(t)=\dot{W}(t)\hfill \\ \mathrm{\Lambda }[Y]=c\hfill \end{array}\text{and}\{\begin{array}{c}DY(t)+A(t)Y(t)=\dot{W}(t)\hfill \\ \mathrm{\Lambda }[Y]=c.\hfill \end{array}$$

Then $`Y^N(t)`$ converges to $`Y(t)`$ pointwise and in $`L^p`$, for all $`p1`$, uniformly in $`t`$, that is:

$$\underset{N+\mathrm{}}{lim}\underset{0t1}{sup}\left|Y^N(t)(\omega )Y(t)(\omega )\right|=0,\omega \mathrm{\Omega },$$

$$\underset{N+\mathrm{}}{lim}\underset{0t1}{sup}Y^N(t)Y(t)_{L^p(\mathrm{\Omega })}=0.$$

Proof: In the situation given, the fundamental solutions $`\mathrm{\Phi }^t(u)^N`$ converge to the fundamental solution $`\mathrm{\Phi }^t(u)`$ uniformly in $`t`$ and $`u`$. From this fact one shows easily that $`J^N`$ and $`(J^N)^1`$ defined by (2.9) converge uniformly to $`J`$ and $`J^1`$, taking into account that the entries of $`dF`$ are finite measures. Hence, the Green functions $`G^N(t,s)`$ converge to $`G(t,s)`$ uniformly in $`t`$ and $`s`$ as well.

We have

$$Y^N(t)Y(t)_{L^p}\left|J^N(t)^1cJ(t)^1c\right|+_0^1\left(G^N(t,s)G(t,s)\right)𝑑B_s_{L^p}.$$

The first term tends to zero uniformly in $`t`$. For the second, note that

$`\underset{t}{sup}\text{E}\left[\left({\displaystyle _0^1}\left(G_{i,j}^N(t,s)G_{i,j}(t,s)\right)𝑑W_s\right)^2\right]=\underset{t}{sup}{\displaystyle _0^1}\left(G_{i,j}^N(t,s)G_{i,j}(t,s)\right)^2𝑑s`$
$`{\displaystyle _0^1}\underset{t}{sup}\left(G_{i,j}^N(t,s)G_{i,j}(t,s)\right)^2ds\underset{t}{sup}\underset{s}{sup}\left(G_{i,j}^N(t,s)G_{i,j}(t,s)\right)^2,`$
which converges to zero. Since all random variables are Gaussian, the convergence to zero of the second moments (uniformly in $`t`$) implies the convergence to zero of all moments, also uniformly in $`t`$. We have proved the second statement of the Theorem.

We turn to the pointwise convergence: Since $`G_{i,j}^N(t,s)G_{i,j}(t,s)`$ is a function of bounded variation which tends to zero uniformly in $`t`$ and $`s`$, the finite measures $`G_{i,j}^N(t,ds)G_{i,j}(t,ds)`$ tend weakly to zero, uniformly in $`t`$, and we have

$$\left|_0^1(G_{i,j}^N(t,s)dW_sG_{i,j}(t,s)dW_s)\right|=\left|_0^1W_s\left(G_{i,j}^N(t,ds)G_{i,j}(t,ds)\right)\right|0.$$

## 3 A characterisation of conditional independence and convergence of conditional expectations

In this section we state two facts of a general nature that will be our main probabilistic tools in Section 4. Lemma 3.4 is an abstract result on the conditional independence of two random vectors when a function of them (of a special structure) is given; it was proved in  (see also ). Proposition 3.7, on the other hand, provides a sufficient condition for the $`L^2`$-convergence as $`N\mathrm{}`$ of a sequence of conditional expectations of the form $`\text{E}[F(U_1^N)|U_2^N].`$

We will mention first three auxiliary lemmas on the conditional independence of $`\sigma `$-fields, whose proofs are not difficult. Recall that we write $`_1\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2`$ to mean that the $`\sigma `$-fields $`_1`$ and $`_2`$ are conditionally independent given the $`\sigma `$-field $`𝒢`$.

###### Lemma 3.1

Let $`_1,_2,𝒢,_1^{},_2^{}`$ be $`\sigma `$-fields such that $`_1^{}_1𝒢`$ and $`_2^{}_2𝒢`$. Then,

$$_1\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2_1^{}\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2^{}.$$

###### Lemma 3.2

Let $`_1,_2,_3`$ and $`𝒢`$ be $`\sigma `$-fields such that $`_1_2\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_3`$ and $`_1\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2`$. Then,

$$_1_3\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2\text{and}_2_3\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_1.$$

###### Lemma 3.3

Let $`_1,_2`$ and $`𝒢`$ be $`\sigma `$-fields such that $`_1\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2`$ and $`𝒢_1`$. Then, for any $`\sigma `$-field $``$, with $`𝒢_1`$,

$$_1\underset{}{\text{ }\text{ }\text{ }\text{ }\text{ }}_2.$$

Let $`(\mathrm{\Omega },,P)`$ be a probability space and $`_1`$ and $`_2`$ two independent sub-$`\sigma `$-fields of $``$. Consider two functions $`g_1:IR^d\times \mathrm{\Omega }IR^d`$ and $`g_2:IR^d\times \mathrm{\Omega }IR^d`$ such that $`g_i`$ is $`(IR^d)_i`$–measurable, $`i=1,2`$. Set $`B(\epsilon ):=\{xIR^d,|x|<\epsilon \}`$, and denote by $`\lambda `$ the Lebesgue measure on $`IR^d`$. Let us introduce the following hypotheses:

There exists $`\epsilon _0>0`$ such that for almost all $`\omega \mathrm{\Omega }`$, and for any $`|\xi |<\epsilon _0`$, $`|\eta |<\epsilon _0`$ the system

$$\{\begin{array}{c}z_1g_1(z_2,\omega )=\xi \hfill \\ z_2g_2(z_1,\omega )=\eta \hfill \end{array}$$

has a unique solution $`(z_1,z_2)IR^{2d}`$.

For every $`z_1IR^d`$ and $`z_2IR^d`$, the random vectors $`g_1(z_2,)`$ and $`g_2(z_1,)`$ possess absolutely continuous distributions and the function

$$\delta (z_1,z_2)=\underset{0<\epsilon <\epsilon _0}{sup}\frac{1}{\lambda (B(\epsilon ))^2}P\{\left|z_1g_1(z_2)\right|<\epsilon ,\left|z_2g_2(z_1)\right|<\epsilon \}$$

is locally integrable in $`IR^{2d}`$, for some $`\epsilon _0>0`$.

For almost all $`\omega \mathrm{\Omega }`$, the functions $`z_2g_1(z_2,\omega )`$ and $`z_1g_2(z_1,\omega )`$ are continuously differentiable and

$$\underset{\genfrac{}{}{0pt}{}{\left|z_2g_2(z_1,\omega )\right|<\epsilon _0}{\left|z_1g_1(z_2,\omega )\right|<\epsilon _0}}{sup}\left|det[\text{I}g_1(z_2,\omega )g_2(z_1,\omega )]\right|^1L^1(\mathrm{\Omega })$$

for some $`\epsilon _0>0`$, where $`g_i`$ denotes the Jacobian matrix of $`g_i`$ with respect to the first argument.

Note that hypothesis (H1) implies the existence of two random vectors $`Z_1`$ and $`Z_2`$ determined by the system

$$\{\begin{array}{c}Z_1(\omega )=g_1(Z_2(\omega ),\omega )\hfill \\ Z_2(\omega )=g_2(Z_1(\omega ),\omega )\hfill \end{array}$$

###### Lemma 3.4

(Alabert, Ferrante, Nualart ). Suppose the functions $`g_1`$ and $`g_2`$ satisfy the above hypotheses (H1) to (H3). Then the following statements are equivalent:

$`_1`$ and $`_2`$ are conditionally independent given the random vectors $`Z_1,Z_2`$.

There exist two functions $`F_i:IR^{2d}\times \mathrm{\Omega }IR`$, $`i=1,2`$, which are $`(IR^{2d})_i`$–measurable, such that

$$\left|det[\text{I}g_1(Z_2)g_2(Z_1)]\right|=F_1(Z_1,Z_2,\omega )F_2(Z_1,Z_2,\omega ),\text{a.s.}$$

Goggin  gives a sufficient condition for the convergence in distribution of a sequence of conditional expectations of the form $`\text{E}[F(U_1^N)|U_2^N]`$. We reproduce here a slightly simplified version. Combining this result with Lemma 3.6, due to Knudsen , we can easily prove our Proposition 3.7.

###### Lemma 3.5

(Goggin ). Let $`U_1^N`$ and $`U_2^N`$ be two sequences of random vectors on a probability space $`(\mathrm{\Omega },,\text{P})`$, such that $`(U_1^N,U_2^N)(U_1,U_2)`$ as $`N\mathrm{}`$ in distribution. Assume that:

1. There exist probabilities $`Q^N`$ on $`(\mathrm{\Omega },)`$ such that $`PQ^N`$ on $`\sigma \{U_1^N,U_2^N\}`$ and $`U_1^N`$ and $`U_2^N`$ are independent under $`Q^N`$. Denote $`\mathrm{}^N(U_1^N,U_2^N):=\frac{dP}{dQ^N}`$ .
2. There exists a probability $`Q`$ on $`(\mathrm{\Omega },)`$ under which $`U_1`$ and $`U_2`$ are independent.
3. The $`Q^N`$–distribution of $`(U_1^N,U_2^N,\mathrm{}^N(U_1^N,U_2^N))`$ converges weakly to the $`Q`$–distribution of $`(U_1,U_2,\mathrm{}(U_1,U_2))`$, where $`\mathrm{}`$ is such that $`\text{E}_Q[\mathrm{}(U_1,U_2)]=1`$.

Then:

1. $`PQ`$ on $`\sigma \{U_1,U_2\}`$ and $`\frac{dP}{dQ}=\mathrm{}(U_1,U_2)`$;
2. For every bounded continuous function $`F`$,

$$\text{E}_P\left[F(U_1^N)|U_2^N\right]E_P\left[F(U_1)|U_2\right]\text{ in distribution.}$$

###### Lemma 3.6

(Knudsen ). Let $`U_1^N`$ and $`U_2^N`$ be two sequences of random vectors on a probability space $`(\mathrm{\Omega },,𝒫)`$. Assume that, as $`N\mathrm{}`$,

1. $`U_2^N\stackrel{P}{}U_2`$.
2. $`U_1^N\stackrel{L^2}{}U_1`$, with $`U_1L^p`$, for some $`p>2`$.
3. $`\text{E}[U_1|U_2^N]_{L^2}\text{E}[U_1|U_2]_{L^2}`$.

Then

$$\text{E}\left[U_1^N|U_2^N\right]\stackrel{L^2}{}\text{E}\left[U_1|U_2\right]$$

as $`N\mathrm{}`$.

Combining Lemma 3.5 and 3.6, we get the following proposition.

###### Proposition 3.7

Let $`U_1^N`$ and $`U_2^N`$ be two sequences of random vectors. Assume that, as $`N\mathrm{}`$,

1. $`U_1^N\stackrel{P}{}U_1`$ and $`U_2^N\stackrel{P}{}U_2`$.
2. Hypotheses 1,2 and 3 of Lemma 3.5 hold true.

Then, for any bounded and continuous function $`F`$,

$$\text{E}\left[F(U_1^N)|U_2^N\right]\stackrel{L^2}{}\text{E}\left[F(U_1)|U_2\right].$$

Proof: Applying Lemma 3.5, we have that, for every bounded and continuous $`F`$,

(3.1)
$$\text{E}\left[F(U_1^N)|U_2^N\right]\text{E}\left[F(U_1)|U_2\right]\text{ in distribution.}$$

From (3.1) and the fact that $`F(U_1^N),F(U_1)L^{\mathrm{}}`$, we obtain the convergence of the $`L^2`$ norms:

(3.2)
$$\text{E}[F(U_1^N)|U_2^N]_{L^2}\text{E}[F(U_1)|U_2]_{L^2}.$$

On the other hand, since $`F`$ is bounded and $`U_1^N\stackrel{P}{}U_1`$, we also have

(3.3)
$$F(U_1^N)\stackrel{L^2}{}F(U_1).$$

Now, (3.2) and (3.3) imply that

$$\text{E}[F(U_1)|U_2^N]_{L^2}\text{E}[F(U_1)|U_2]_{L^2}$$

and we get the conclusion applying Lemma 3.6.

## 4 Markovian properties of linear functional boundary value problems

In the study of boundary value stochastic problems, one of the main interests has been to seek conditions on the coefficients for the solution process to satisfy some suitably defined Markov-type property. Intuition suggests that a relation $`h(X(0),X(1))=0`$ will possibly prevent the Markov process property from holding in general. One might think that nevertheless the Markov field property, which is defined below, will be satisfied. It is easy to see that any Markov process is a Markov field (see Jamison  for the continuous case and Alabert and Marmolejo  for a simple proof in the general case). The converse is not true. For instance, the processes $`X(t)=W(t)\alpha W(1)`$ are Markov fields; they are not Markov processes, except for the cases $`\alpha =0`$ and $`\alpha =1`$.

###### Definition 4.1

A process $`\{X(t),t[0,1]\}`$ is said to be a Markov field if for any $`0a<b1`$, the $`\sigma `$-fields $`\sigma \{X(t),t[a,b]\}`$ and $`\sigma \{X(t),t]a,b[^c\}`$ are conditionally independent given $`\sigma \{X(a),X(b)\}`$.

However, even this weaker property holds only in special cases. For instance, in  it was shown that the solution to

(4.1)
$$\{\begin{array}{c}\dot{X}(t)=b(X(t))+\sigma (X(t))\dot{W}(t),t[0,1]\hfill \\ X(0)=\psi (X(1)),\hfill \end{array}$$

where the stochastic integral is understood in the Stratonovich sense, is a Markov field if and only if $`b(x)=A\sigma (x)+B\sigma (x){\displaystyle _c^x}{\displaystyle \frac{1}{\sigma (t)}}𝑑t`$, for some constants $`A,B,c`$. As a corollary, in case $`\sigma `$ is a constant (additive noise), $`X`$ is a Markov field if and only if $`b`$ is an affine function.

Our aim is to study the linear–additive case when the additional condition takes into account the value of the solution in some interior points of the time interval. The following simple example illustrates that the situation changes.

###### Example 4.2

Consider the first order system

$$\{\begin{array}{c}\dot{X}(t)=\dot{W}(t),t[0,1]\hfill \\ X(\frac{1}{2})+X(1)=0.\hfill \end{array}$$

The solution is the process

$$X(t)=\frac{1}{2}\left(W(\frac{1}{2})+W(1)\right)+W(t),$$

which is not a Markov field. Indeed, for $`a=0`$ and $`b=\frac{2}{3}`$, the random variables $`X(\frac{1}{2})`$ and $`X(1)`$ are not conditionally independent given $`\sigma \{X(a),X(b)\}`$. Nevertheless, $`X`$ is a Markov field when restricted to $`[0,\frac{1}{2}]`$ or $`[\frac{1}{2},1]`$.

In order to formulate precisely the conditional independence property enjoyed by the system (2.3)-(2.4), we introduce first some more concepts and notation.

Let $`\mathrm{\Lambda }_1,\mathrm{},\mathrm{\Lambda }_n`$ be the real-valued components of a boundary operator $`\mathrm{\Lambda }`$ of the form (2.4) and denote their support by $`\text{supp}\mathrm{\Lambda }_i:=\{t_j[0,1]:\alpha _{ij}0\}`$.

###### Definition 4.3

We will say that $`\mathrm{\Lambda }_i`$ preserves the pair $`(a,b)`$ if either $`\text{supp}\mathrm{\Lambda }_i]a,b[`$ or $`\text{supp}\mathrm{\Lambda }_i[a,b]^c`$. If this is true for all $`i`$ (i.e. there are no boundary conditions involving simultaneously points inside and outside $`[a,b]`$), then we will also say that $`\mathrm{\Lambda }`$ preserves $`(a,b)`$.

We want to prove that the solution $`Y`$ to the system (2.3)-(2.4) satisfies the following conditional independence property (Theorem 4.5): If $`\mathrm{\Lambda }`$ preserves $`(a,b)`$, then

$$\sigma \{Y(t),t[a,b]\}\underset{\sigma \{Y(a),Y(b)\}}{\text{ }\text{ }\text{ }\text{ }\text{ }}\sigma \{Y(t),t]a,b[^c\}.$$

More generally, this conditional independence is also true when $`\mathrm{\Lambda }`$ does not preserve $`(a,b)`$, provided the conditioning $`\sigma `$-field is enlarged with the variables $`Y_n(t)`$, for $`t`$ in $`[a,b]`$ and in the support of all non-preserving boundary operators $`\mathrm{\Lambda }_i`$ (Theorem 4.10).

Since the boundary conditions can be written in many different equivalent ways, and the sets $`\text{supp}\mathrm{\Lambda }_i`$ (hence the property of preserving an interval) depend on the representation chosen, we need, before proceeding further, some sort of “canonical” definition of the linear operator $`\mathrm{\Lambda }`$. Given $`\text{supp}\mathrm{\Lambda }=\{t_1,\mathrm{},t_m\}`$, $`\mathrm{\Lambda }`$ can be regarded as a linear mapping $`IR^mIR^n`$, that means, an $`n\times m`$ matrix acting on the vector $`(Y_n(t_1),\mathrm{},Y_n(t_m))`$ (see notations of Section 2).

A basis $`B`$ for an $`n\times m`$ matrix $`\mathrm{\Lambda }`$ is any $`n\times n`$ minor with full rank. For notational simplicity, assume that $`B`$ consists of the firsts $`n`$ columns of $`\mathrm{\Lambda }`$. Denoting by $`N`$ the non-basic columns, we can write $`\mathrm{\Lambda }=(B,N)`$. Defining $`\stackrel{~}{\mathrm{\Lambda }}=(I,B^1N)`$, the system of equations $`\mathrm{\Lambda }x=c`$ can be written in the equivalent form $`\stackrel{~}{\mathrm{\Lambda }}x=B^1c`$. In this situation, we shall say that $`\stackrel{~}{\mathrm{\Lambda }}`$ is a basic expression of $`\mathrm{\Lambda }`$ relative to the basis $`B`$. In the following lemma we prove that this representation can be considered “canonical” for our purposes, since any pair $`(a,b)`$ will or will not be preserved by any basic equivalent form of $`\mathrm{\Lambda }`$. In the sequel, we will always assume, without explicit mention, that the boundary condition is written in this form.

###### Lemma 4.4

Let $`\stackrel{~}{\mathrm{\Lambda }}`$ and $`\stackrel{~}{\mathrm{\Lambda }}^{}`$ be two basic expressions of $`\mathrm{\Lambda }`$, and fix $`0a<b1`$. Then, $`\stackrel{~}{\mathrm{\Lambda }}`$ preserves $`(a,b)`$ if and only if $`\stackrel{~}{\mathrm{\Lambda }}^{}`$ preserves $`(a,b)`$.

Proof: Without any loss of generality we can assume that $`\stackrel{~}{\mathrm{\Lambda }}=(\text{I},N)`$, where I is the $`n\times n`$ identity matrix and

$$N=\left[\begin{array}{ccc}\alpha _{1,n+1}& \mathrm{}& \alpha _{1,m}\\ \mathrm{}& \mathrm{}& \mathrm{}\\ \alpha _{n,n+1}& \mathrm{}& \alpha _{n,m}\end{array}\right].$$

All basic expressions of the original matrix $`\mathrm{\Lambda }`$ can be obtained by repeated Gaussian pivoting on entries of non-basic columns; when pivoting on $`\alpha _{ik}`$, the column $`i`$ leaves the basis (the identity matrix) and is replaced by column $`k`$. Therefore, it is sufficient to prove the lemma for $`\stackrel{~}{\mathrm{\Lambda }}`$ and a basic expression $`\stackrel{~}{\mathrm{\Lambda }}^{}`$ obtained from $`\stackrel{~}{\mathrm{\Lambda }}`$ by one pivoting operation.

Let us assume that $`\alpha _{1,n+1}0`$ and that the operator $`\stackrel{~}{\mathrm{\Lambda }}_i`$ does not preserve the pair $`(a,b)`$. We are going to find an operator $`\stackrel{~}{\mathrm{\Lambda }}_j^{}`$ which neither preserves $`(a,b)`$. Rows $`1`$ and $`i`$ before and after pivoting on $`\alpha _{1,n+1}`$ are the following:

$$\begin{array}{ccccccccccccccc}& & & & & & i& & & & & & & & \\ \stackrel{~}{\mathrm{\Lambda }}_1=& [& 1& 0& \mathrm{}& 0& 0& 0& \mathrm{}& 0& \alpha _{1,n+1}& \alpha _{1,n+2}& \mathrm{}& \alpha _{1,m}& ]\\ \stackrel{~}{\mathrm{\Lambda }}_i=& [& 0& 0& \mathrm{}& 0& 1& 0& \mathrm{}& 0& \alpha _{i,n+1}& \alpha _{i,n+2}& \mathrm{}& \alpha _{i,m}& ]\\ \stackrel{~}{\mathrm{\Lambda }}_1^{}=& [& \frac{1}{\alpha _{1,n+1}}& 0& \mathrm{}& 0& 0& 0& \mathrm{}& 0& 1& \frac{\alpha _{1,n+2}}{\alpha _{1,n+1}}& \mathrm{}& \frac{\alpha _{1,m}}{\alpha _{1,n+1}}& ]\\ \stackrel{~}{\mathrm{\Lambda }}_i^{}=& [& \beta _i& 0& \mathrm{}& 0& 1& 0& \mathrm{}& 0& 0& \gamma _{i,n+2}& \mathrm{}& \gamma _{i,m}& ]\end{array}$$

where $`\beta _i=\frac{\alpha _{i,n+1}}{\alpha _{1,n+1}}`$ and $`\gamma _{i,j}=\alpha _{i,j}\alpha _{1,j}\beta _i`$.

If $`\stackrel{~}{\mathrm{\Lambda }}_1`$ does not preserve $`(a,b)`$, then the result is trivially true, with $`j=1`$, since $`\stackrel{~}{\mathrm{\Lambda }}_1`$ and $`\stackrel{~}{\mathrm{\Lambda }}_1^{}`$ have their non-zero coefficients in the same columns. By a similar reason, if $`\beta _i=0`$ we can take $`j=i`$.

Assume finally that $`\text{supp}\stackrel{~}{\mathrm{\Lambda }}_1]a,b[`$ (so that $`i1`$), and that $`\beta _i0`$. If $`t_i[a,b]^c`$, the result is proved, since $`\stackrel{~}{\mathrm{\Lambda }}_i^{}`$ links $`t_1`$ and $`t_i`$, and $`t_1]a,b[`$; take $`j=i`$. If $`t_i]a,b[`$, then there exists $`k\{n+1,\mathrm{},m\}`$ such that $`t_k[a,b]^c`$ with $`\alpha _{ik}0`$; since $`\alpha _{1k}=0`$ (because $`\text{supp}\stackrel{~}{\mathrm{\Lambda }}_1]a,b[`$), we have $`\gamma _{ik}=\alpha _{ik}0`$ and we can take again $`j=i`$.

Now we can formulate our main result:

###### Theorem 4.5

Suppose the system

$$\{\begin{array}{c}DY(t)+A(t)Y(t)=\dot{B}(t),t[0,1],\hfill \\ \mathrm{\Lambda }[Y]=c\hfill \end{array}$$

satisfies (H0), and let $`Y=\{Y(t),t[0,1]\}`$ be its unique solution. Then,

(4.2)
$$\sigma \{Y(t),t[a,b]\}\underset{\sigma \{Y(a),Y(b)\}}{\text{ }\text{ }\text{ }\text{ }\text{ }}\sigma \{Y(t),t]a,b[^c\}$$

if and only if the pair $`(a,b)`$ is preserved by $`\mathrm{\Lambda }`$.

Our main tool for the proof of the ‘if’ part in Theorem 4.5 will be Lemma 3.4. The idea is the following: We will split the $`2n`$-dimensional random vector $`(Y(a),Y(b))`$ into two vectors $`Z^1`$ and $`Z^2`$ of suitable dimensions, in such a way that $`Z^1`$ be a function of $`Z^2`$ and the increments of the Wiener process $`W`$ in $`[a,b]`$, and in turn $`Z^2`$ be a function of $`Z^1`$ and the increments of $`W`$ in $`]a,b[^c`$. These mappings will play the role of $`g_1`$ and $`g_2`$ in the set of hypotheses (H1) to (H3). The first will be defined through the solution to equation $`DY(t)+A(t)Y(t)=\dot{B}(t)`$, with the components of $`Z^2`$ fixed to a constant; the second will be defined similarly, fixing the components of $`Z^1`$ to a constant. However, this means that we need to solve our differential equation with several sets of constraints, which are different from the original set, and therefore we cannot ensure a priori that these problems are well-posed. Consequently, the above functions $`g_1`$ and $`g_2`$ need not exist in general.

To solve this technical difficulty, we will resort to a two-step procedure. First, we will assume that all functional boundary value problems that we need to solve are indeed well-posed. Then, hypotheses (H1) to (H3) can be checked, and Lemma 3.4 applies directly, yielding the desired result. This is the goal of Proposition 4.7. Secondly, we will use the approximation result of Subsection 2.3 to show that the matrix $`A`$ can be approximated by perturbed matrices $`A^N`$ for which all boundary problems involved are well-posed and whose solutions $`Y^N`$ converge to the solution $`Y`$ of the original problem. Then, the convergence of conditional expectations given in Proposition 3.7 will allow to carry the conditional independence properties of $`Y^N`$ to the limit. This second step is the contents of Proposition 4.8. The ‘only if’ part of the theorem is shown in Proposition 4.9.

Let us formulate precisely the assumption needed for the first step: Set

(4.3)
$$\begin{array}{c}\mathrm{}=\mathrm{\#}\{i:\text{supp}\mathrm{\Lambda }_i[0,a[\},\hfill \\ p=\mathrm{\#}\{i:\text{supp}\mathrm{\Lambda }_i]b,1]\},\hfill \\ q=\mathrm{\#}\{i:\text{supp}\mathrm{\Lambda }_i]a,b[\}.\hfill \end{array}$$

We can assume that the equalities $`\mathrm{\Lambda }_i[X]=c_i`$ are ordered in the following way:

(4.4)
$$\begin{array}{c}\text{supp}\mathrm{\Lambda }_i[0,a[,i=1,\mathrm{},\mathrm{},\hfill \\ \text{supp}\mathrm{\Lambda }_i]a,b[,i=\mathrm{}+1,\mathrm{},\mathrm{}+q,\hfill \\ \text{supp}\mathrm{\Lambda }_i]b,1],i=\mathrm{}+q+1,\mathrm{},\mathrm{}+q+p,\hfill \end{array}$$

and the remaining equations (those involving points both in $`[0,a[`$ and in $`]b,1]`$), carry the labels $`i=\mathrm{}+q+p+1,\mathrm{},n`$.

Consider now $`DY(t)+A(t)Y(t)=0`$ with the following sets of lateral conditions and the specified domain:

(4.5)
$$\{\begin{array}{c}Y_j(a)=0,j=1,\mathrm{},n\mathrm{}\hfill \\ \mathrm{\Lambda }_i[Y]=0,i=1,\mathrm{},\mathrm{}\hfill \end{array},\text{on}[0,a];$$

(4.6)
$$\{\begin{array}{c}Y_j(b)=0,j=1,\mathrm{},n\mathrm{}q\hfill \\ Y_j(a)=0,j=n\mathrm{}+1,\mathrm{},n\hfill \\ \mathrm{\Lambda }_i[Y]=0,i=\mathrm{}+1,\mathrm{},\mathrm{}+q\hfill \end{array},\text{on}[a,b];$$

(4.7)
$$\{\begin{array}{c}Y_j(b)=0,j=n\mathrm{}q+1,\mathrm{},n\hfill \\ \mathrm{\Lambda }_i[Y]=0,i=\mathrm{}+q+1,\mathrm{},\mathrm{}+q+p\hfill \\ \mathrm{\Lambda }_i[Y]=0,i=\mathrm{}+q+p+1,\mathrm{},n\hfill \\ Y_n(t)=0,t\underset{i=\mathrm{}+q+p+1}{\overset{n}{}}\left(\text{supp}\mathrm{\Lambda }_i[0,a]\right)\hfill \end{array},\text{on}[b,1],$$

(notice that the third and fourth lines result in $`n\mathrm{}qp`$ equations involving only points in $`[b,1]`$).

###### Definition 4.6

We will say that the pair $`(a,b)`$ is regular if $`DY(t)+A(t)Y(t)=0`$ together with any of the sets of conditions (4.5), (4.6) or (4.7) has only the trivial solution. Otherwise $`(a,b)`$ will be called singular.

###### Proposition 4.7

Suppose the system

(4.8)
$$\{\begin{array}{c}DY(t)+A(t)Y(t)=\dot{B}(t),t[0,1],\hfill \\ \mathrm{\Lambda }[Y]=c\hfill \end{array}$$

satisfies (H0), and let $`Y=\{Y(t),t[0,1]\}`$ be its unique solution. Let $`(a,b)`$ be a regular pair preserved by $`\mathrm{\Lambda }`$. Then (4.2) holds true.

Proof: Let us define the $`\sigma `$-fields

$$\begin{array}{c}_{a,b}^i=\sigma \{W_tW_a,t[a,b]\}\hfill \\ _{a,b}^e=\sigma \{W_t,t[0,a]\}\sigma \{W_1W_t,t[b,1]\}\hfill \end{array}$$

for $`0a<b1`$. Notice that $`_{a,b}^i`$ and $`_{a,b}^e`$ are independent.

We shall divide the proof into several steps. In Step 1 we reduce the proof to that of the conditional independence of two independent $`\sigma `$-fields. In Step 2 it is shown that there exist the two functions $`g_1`$ and $`g_2`$ needed to apply Lemma 3.4. The hypotheses of this lemma are checked in Steps 3, 4 and 5. In Step 6 we finally conclude the result.

Step 1 Denote $`𝒢_{a,b}=\sigma \{Y(a),Y(b)\}`$. If

(4.9)
$$_{a,b}^i\underset{𝒢_{a,b}}{\text{ }\text{ }\text{ }\text{ }\text{ }}_{a,b}^e,$$

then (4.2) holds.

Proof of Step 1: It is immediate to prove that $`\sigma \{Y(t),t[a,b]\}𝒢_{a,b}_{a,b}^i`$, and $`\sigma \{Y(t),t]a,b[^c\}𝒢_{a,b}_{a,b}^e`$. We apply then Lemma 3.1.

Step 2 Let $`\mathrm{}`$, $`p`$ and $`q`$ be as in $`(\text{4.3})`$. We will denote by $`\stackrel{~}{Y}`$ the solution to (4.8), to distinguish the actual solution from $`Y`$ regarded as an unknown of the system. Define

$$\begin{array}{c}Z^1:=(\stackrel{~}{Y}_1(a),\mathrm{},\stackrel{~}{Y}_n\mathrm{}(a),\stackrel{~}{Y}_{n\mathrm{}q+1}(b),\mathrm{},\stackrel{~}{Y}_n(b))IR^{n+q},\hfill \\ Z^2:=(\stackrel{~}{Y}_1(b),\mathrm{},\stackrel{~}{Y}_{n\mathrm{}q}(b),\stackrel{~}{Y}_{n\mathrm{}+1}(a),\mathrm{},\stackrel{~}{Y}_n(a))IR^{nq}.\hfill \end{array}$$

Then, there exist two functions

$$\begin{array}{c}g_1:IR^{nq}\times \mathrm{\Omega }IR^{n+q},\hfill \\ g_2:IR^{n+q}\times \mathrm{\Omega }IR^{nq},\hfill \end{array}$$

measurable with respect to $`(IR^{nq})_{a,b}^i`$ and $`(IR^{n+q})_{a,b}^e`$ respectively, and such that

(4.10)
$$Z^1=g_1(Z^2,\omega )\text{and}Z^2=g_2(Z^1,\omega ).$$

Proof of Step 2: Consider the lateral conditions

(4.11)
$$\{\begin{array}{c}Y_j(b)=Z_j^2,j=1,\mathrm{},n\mathrm{}q\hfill \\ Y_j(a)=Z_{jq}^2,j=n\mathrm{}+1,\mathrm{},n\hfill \\ \mathrm{\Lambda }_i[Y]=c_i,i=\mathrm{}+1,\mathrm{},\mathrm{}+q\hfill \end{array}$$

on $`[a,b]`$. The process $`\stackrel{~}{Y}`$ trivially satisfies $`DY(t)+A(t)Y(t)=\dot{B}(t)`$ and these conditions on $`[a,b]`$; however the solution to this problem is also unique. Therefore, taking into account that $`\mathrm{\Lambda }_i[\stackrel{~}{Y}]=c_i`$ are constants, the vector $`Z^1`$ is determined by $`Z^2`$ and the increments of the Wiener process in $`[a,b]`$. Moreover, the function $`g_1(z_2,\omega )`$ so defined has a sense for every $`z_2IR^{nq}`$, because we have that the solution to (2.3)-(4.11) is unique, and this fact does not depend on the particular right-hand sides.

We want to prove analogously the existence of the function $`g_2`$. Consider first $`DY(t)+A(t)Y(t)=\dot{B}(t)`$ on $`[0,a]`$ with conditions

(4.12)
$$\{\begin{array}{c}Y_j(a)=Z_j^1,j=1,\mathrm{},n\mathrm{}\hfill \\ \mathrm{\Lambda }_i[Y]=c_i,i=1,\mathrm{},\mathrm{}.\hfill \end{array}$$

The restriction to $`[0,a]`$ of the solution $`\stackrel{~}{Y}`$ to (4.8) solves also the differential system with conditions (4.12), and is its unique solution. Consider now $`DY(t)+A(t)Y(t)=\dot{B}(t)`$ on $`[b,1]`$, with

(4.13)
$$\{\begin{array}{c}Y_j(b)=Z_{j+q}^1,j=n\mathrm{}q+1,\mathrm{},n\hfill \\ \mathrm{\Lambda }_i[Y]=c_i,i=\mathrm{}+q+1,\mathrm{},\mathrm{}+q+p\hfill \end{array}$$

and the $`n\mathrm{}qp`$ equations on $`[b,1]`$ that result from

(4.14)
$$\{\begin{array}{c}\mathrm{\Lambda }_i[Y]=c_i,i=\mathrm{}+q+p+1,\mathrm{},n\hfill \\ Y_n(t)=\stackrel{~}{Y}_n(t),t\underset{i=\mathrm{}+q+p+1}{\overset{n}{}}\left(\text{supp}\mathrm{\Lambda }_i[0,a]\right).\hfill \end{array}$$

Again, $`\stackrel{~}{Y}`$ restricted to $`[b,1]`$ is its unique solution with conditions (4.13)-(4.14). The values $`\stackrel{~}{Y}_n(t)`$ appearing here are found in (4.13) as a function of $`Z_{j+q}^1`$, $`j=n\mathrm{}q+1,\mathrm{},n`$, and the Wiener process on $`[0,a]`$. Therefore, the whole vector $`Z^2`$ is determined by $`Z^1`$ and the increments of $`W`$ in $`]a,b[^c`$. As before, the function $`g_2(z_1,\omega )`$ so defined has a sense for all $`z_1IR^{n+q}`$.

Step 3 The functions $`g_1`$ and $`g_2`$ found in Step 2 satisfy (H1).

Proof of Step 3: The solution to a linear differential equation depends linearly on the lateral data $`c`$ (see (2.11)). Therefore, for each $`\omega `$ fixed, system (4.10) is linear and it is enough to check that it has a unique solution for $`\xi =\eta =0`$.

Now, gathering together the lateral conditions (4.12), (4.13), (4.14), we obtain the original lateral conditions, so that system (4.10) is equivalent to (4.8) and therefore the solution exists and is unique.

Step 4 $`g_1`$ and $`g_2`$ satisfy (H2).

Proof of Step 4: The boundary value problem that defines $`g_1`$ consists of equation $`DY(t)+A(t)Y(t)=\dot{B}(t)`$, together with the conditions (4.11). The resulting vector

$$(Y_1(a),\mathrm{},Y_n\mathrm{}(a),Y_{n\mathrm{}q+1}(b),Y_n(b))$$

is absolutely continuous on $`IR^{n+q}`$, by Proposition 2.1. The proof for $`g_2`$ is analogous, using Remark 2.2.

Finally, the random vectors $`z_1g_1(z_2,\omega )`$ and $`z_2g_2(z_1,\omega )`$ are independent and have the form $`z_1M^1z_2+U^1(\omega )`$ and $`z_2M^2z_1+U^2(\omega )`$ respectively, for some constant matrices $`M^1`$ and $`M^2`$ and some Gaussian absolutely continuous vectors $`U^1`$ and $`U^2`$. We deduce that the $`IR^{2n}`$-valued random vector $`(z_1g_1(z_2,\omega ),z_2g_2(z_1,\omega ))`$ has a density which is uniformly bounded in $`z_1`$ and $`z_2`$. It follows at once that the function $`\delta `$ in (H2) is bounded.

Step 5 $`g_1`$ and $`g_2`$ satisfy (H3). Specifically, $`det[\text{I}g_1(z_2,\omega )g_2(z_1,\omega )]`$ is a constant different from zero.

Proof of Step 5: $`g_1`$ and $`g_2`$ are affine functions of the first argument, with a non-random linear coefficient (see (2.11)). Therefore, $`g_1(z_2,\omega )`$ and $`g_2(z_1,\omega )`$ are constant matrices of dimensions $`(n+q)\times (nq)`$ and $`(nq)\times (n+q)`$ respectively, which we denote simply $`g_1`$ and $`g_2`$. We know that the linear system

$$\{\begin{array}{c}z_1=g_1(z_2,\omega )\hfill \\ z_2=g_2(z_1,\omega )\hfill \end{array}$$

admits a unique solution. This is equivalent to say

$$det[\text{I}g_1g_2]=det\left(\begin{array}{cc}\text{I}& g_1\\ g_2& \text{I}\end{array}\right)0.$$

Step 6 Relation (4.2) holds true.

Proof of Step 6: We can apply Lemma 3.4 and the factorization in (ii) trivially holds. We deduce the relation (4.9) and, by Step 1, that the process $`Y`$ satisfies the desired property, for $`(a,b)`$ regular.

Let us now extend Proposition 4.7 to singular pairs $`(a,b)`$, using an approximation argument. We denote by Eq($`A,\mathrm{\Lambda }`$) our functional boundary value problem relative to the matrix function $`A`$ and the boundary operator $`\mathrm{\Lambda }`$. The boundary data $`c`$ will be fixed throughout. Let us call $`\mathrm{\Lambda }_1`$, $`\mathrm{\Lambda }_2`$, $`\mathrm{\Lambda }_3`$ the operators associated to the lateral conditions given by (4.5), (4.6) and (4.7), respectively.

###### Proposition 4.8

Proposition 4.7 holds also for singular pairs $`(a,b)`$.

Proof: Our initial hypothesis (H0) states that the original problem has one and only one solution, that is, $`A(t)V_\mathrm{\Lambda }`$, where $`V_\mathrm{\Lambda }`$ stands for the set $`V`$ defined in Subsection 2.3, relative to the boundary operator $`\mathrm{\Lambda }`$.

Fix $`0ab`$. If the pair $`(a,b)`$ is singular, then at least one of the problems Eq($`A,\mathrm{\Lambda }_1`$), Eq($`A,\mathrm{\Lambda }_2`$), or Eq($`A,\mathrm{\Lambda }_3`$) is not well-posed. We know from Lemma 2.3 that $`V_\mathrm{\Lambda }`$ is open and dense in the space $`𝒞`$ of matrices of the form (2.5), for any $`\mathrm{\Lambda }`$. Therefore the set $`V:=V_\mathrm{\Lambda }V_{\mathrm{\Lambda }_1}V_{\mathrm{\Lambda }_2}V_{\mathrm{\Lambda }_3}`$ is also open and dense in $`𝒞`$.

Let $`\{A^N(t),NIN\}`$ be a sequence of elements of $`V`$ converging to $`A(t)𝒞`$. From Proposition 2.4, the corresponding solutions $`Y^N(t)`$ converge to $`Y(t)`$ in $`L^p`$.

Fix $`s[a,b]^c`$, and $`r_1,\mathrm{},r_k(a,b)`$. Consider the space $`M`$ defined in (2.15), based on the coordinates $`\{s,r_1,\mathrm{},r_k\}\text{supp}\mathrm{\Lambda }`$. Let $`M^{}`$ be the projection of $`M`$ onto the coordinates $`r_1,\mathrm{},r_k`$. Assume that $`Y^N(s)`$ are non-degenerate (if they are, $`Y(s)`$ will also be a constant, and there is nothing to prove). Using Proposition 2.1, the vector $`U^N:=(Y^N(s),Y^N(r_1),\mathrm{},Y^N(r_k))`$ is a Gaussian vector with some density $`f^N`$ with respect to the Hausdorff measure on $`IR\times M^{}`$. Then there clearly exists an equivalent Gaussian probability with density $`f_0`$ on $`IR\times M^{}`$ whose first coordinate is incorrelated with the remaining ones. Define the probability $`Q^N`$ on $`\mathrm{\Omega }`$ by $`dP=\mathrm{}^NdQ^N`$, with $`\mathrm{}^N=\frac{f^N}{f_0}(U^N)`$. Then hypothesis 1 of Lemma 3.5 is clearly satisfied with $`U_1^N=Y^N(s)`$, $`U_2^N=(Y^N(r_1),\mathrm{},Y^N(r_k))`$. Analogously, one can define a probability $`Q`$ by $`dP=\mathrm{}dQ`$, where $`\mathrm{}=\frac{f}{f_0}(U)`$ and $`f`$ is the density of $`U:=(Y(s),Y(r_1),\mathrm{},Y(r_k))`$, satisfying hypothesis 2 of Lemma 3.5 with $`U_1=Y(s)`$ and $`U_2=(Y(r_1),\mathrm{},Y(r_k))`$.

We prove now that hypothesis 3 also holds true: we want to see that if $`h`$ is a bounded and continuous function, then

(4.15)
$$\underset{N\mathrm{}}{lim}\text{E}_{Q^N}\left[h(U^N,\mathrm{}^N(U^N))\right]=\text{E}_Q\left[h(U,\mathrm{}(U))\right].$$

Notice that

$$\text{E}_{Q^N}\left[h(U^N,\mathrm{}^N(U^N))\right]=_Mh(u,\mathrm{}^N(u))f_0(u)(du),$$

where $``$ denotes the Hausdorff measure, and analogously

$$\text{E}_Q\left[h(U,\mathrm{}(U))\right]=_Mh(u,\mathrm{}(u))f_0(u)(du).$$

Convergence (4.15) is then easily derived through the dominated convergence theorem. We can therefore apply Proposition 3.7 to obtain that

$`L^2\text{}\underset{N}{lim}\text{E}[F(Y^N(s))|Y^N(a),Y^N(r_1),\mathrm{},Y^N(r_n),Y^N(b)]`$
$`=\text{E}[F(Y(s))|Y(a),Y(r_1),\mathrm{},Y(r_n),Y(b)].`$
Analogously, one obtains the limit

$$L^2\text{}\underset{N}{lim}\text{E}[F(Y^N(s))|Y^N(a),Y^N(b)]=\text{E}[F(Y(s))|Y(a),Y(b)].$$

We conclude that the conditional independence property can be carried to the limit and this finishes the proof of the present proposition and consequently of the necessity in Theorem 4.5.

The ‘only if’ part of Theorem 4.5 is far easier to prove:

###### Proposition 4.9

If $`\mathrm{\Lambda }`$ does not preserve $`(a,b)`$, then relation (4.2) is false.

Proof: Let us assume, to keep notations simple, that there is only one boundary operator $`\mathrm{\Lambda }_1`$, which corresponds to a first order equation (the general case can be stated similarly). Assume that $`\mathrm{\Lambda }_1`$ does not preserve $`(a,b)`$ and that (4.2) holds true. The corresponding boundary condition can thus be written as

(4.16)
$$\underset{t_j[a,b]^c}{}\alpha _jY(t_j)+\underset{t_j]a,b[}{}\alpha _jY(t_j)=c,$$

where none of the summations is void.

Notice first that

(4.17)
$$(\underset{t_j]a,b[}{}\alpha _jY(t_j),Y(a),Y(b))$$

is an absolutely continuous random vector in $`IR^3`$. This follows easily from Proposition 2.1.

Now, for any bounded and measurable function $`\mathrm{\Psi }:IRIR`$, we have, by the conditional independence hypothesis and relation (4.16), that

$$\begin{array}{c}\text{E}[\mathrm{\Psi }(_{t_j]a,b[}\alpha _jY(t_j))|Y(a),Y(b)]=\hfill \\ \\ =\text{E}[\mathrm{\Psi }(_{t_j]a,b[}\alpha _jY(t_j))|\{Y(t),t]a,b[^c\}]=\mathrm{\Psi }(_{t_j]a,b[}\alpha _jY(t_j)).\hfill \end{array}$$

Taking $`\mathrm{\Psi }=\text{1}_{[M,M]}`$ for some $`M>0`$, we get that, on $`\{|_{t_j]a,b[}\alpha _jY(t_j)|M\}`$, which is a set of positive probability, $`_{t_j]a,b[}\alpha _jY(t_j)`$ is a measurable function of $`(Y(a),Y(b))`$. In particular, this contradicts the absolute continuity of (4.17).

Next theorem generalises Theorem 4.5 by allowing the existence of non-preserving boundary operators, at the price of enlarging the conditioning $`\sigma `$-field. The result can hardly be called a Markovian type property; nevertheless, it seems interesting in itself, and gives rise to the conjecture contained in Remark 4.12 below.

###### Theorem 4.10

Suppose the system

$$\{\begin{array}{c}DY(t)+A(t)Y(t)=\dot{B}(t),t[0,1],\hfill \\ \mathrm{\Lambda }[Y]=c\hfill \end{array}$$

satisfies (H0), and let $`Y=\{Y(t),t[0,1]\}`$ be its unique solution. Fix $`0a<b1`$ such that $`a,b\text{supp}\mathrm{\Lambda }`$, and let $`𝒢`$ be the $`\sigma `$-field generated by $`Y(a)`$, $`Y(b)`$ and all variables $`Y_n(t)`$, with $`t`$ in $`\text{supp}\mathrm{\Lambda }_i]a,b[`$ for some $`\mathrm{\Lambda }_i`$ not preserving $`(a,b)`$. Then,

(4.18)
$$\sigma \{Y(t),t[a,b]\}\underset{𝒢}{\text{ }\text{ }\text{ }\text{ }\text{ }}\sigma \{Y(t),t]a,b[^c\}.$$

The same holds true replacing $`]a,b[`$ by $`[a,b]^c`$ in the definition of $`𝒢`$.

Proof: The result can be proved as Theorem 4.5 with some modifications. We will only give a sketch of the necessary changes in the simple case where there is only one non-preserving boundary condition, which links one only point $`t^{}`$ inside $`[a,b]`$ with one or more points in $`[0,a]`$. Specifically, fix $`a,b`$ and assume that $`\mathrm{}`$, $`p`$ and $`q`$ are as in (4.3) and (4.4), the boundary conditions with support in $`[0,a][b,1]`$ carry the labels $`i=\mathrm{}+q+p+1,\mathrm{},n1`$, and that $`\mathrm{\Lambda }_n`$ is the non-preserving condition.

In Step 2, define

$$\begin{array}{c}Z^1:=(\stackrel{~}{Y}_1(a),\mathrm{},\stackrel{~}{Y}_{n\mathrm{}1}(a),\stackrel{~}{Y}_{n\mathrm{}q}(b),\mathrm{},\stackrel{~}{Y}_n(b),\stackrel{~}{Y}_n(t^{}))IR^{n+q+1},\hfill \\ Z^2:=(\stackrel{~}{Y}_1(b),\mathrm{},\stackrel{~}{Y}_{n\mathrm{}q1}(b),\stackrel{~}{Y}_n\mathrm{}(a),\mathrm{},\stackrel{~}{Y}_n(a))IR^{nq}.\hfill \end{array}$$

Consider the lateral conditions

$$\{\begin{array}{c}Y_j(b)=Z_j^2,j=1,\mathrm{},n\mathrm{}q1\hfill \\ Y_j(a)=Z_{jq}^2,j=n\mathrm{},\mathrm{},n\hfill \\ \mathrm{\Lambda }_i[Y]=c_i,i=\mathrm{}+1,\mathrm{},\mathrm{}+q\hfill \end{array}$$

on $`[a,b]`$. This system, as before, defines the function $`g_1`$. To define function $`g_2`$, we consider first $`DY(t)+A(t)Y(t)=\dot{B}(t)`$ on $`[0,a]`$ with conditions

$$\{\begin{array}{c}Y_j(a)=Z_j^1,j=1,\mathrm{},n\mathrm{}1\hfill \\ \mathrm{\Lambda }_i[Y]=c_i,i=1,\mathrm{},\mathrm{}\hfill \\ \mathrm{\Lambda }_n[Y]=c_n\hfill \\ Y_n(t^{})=\stackrel{~}{Y}_n(t^{}),\hfill \end{array}$$

and secondly the system on $`[b,1]`$, with

$$\{\begin{array}{c}Y_j(b)=Z_{j+q}^1,j=n\mathrm{}q,\mathrm{},n\hfill \\ \mathrm{\Lambda }_i[Y]=c_i,i=\mathrm{}+q+1,\mathrm{},\mathrm{}+q+p\hfill \end{array}$$

and the $`n1\mathrm{}qp`$ equations on $`[b,1]`$ that result from

$$\{\begin{array}{c}\mathrm{\Lambda }_i[Y]=c_i,i=\mathrm{}+q+p+1,\mathrm{},n1\hfill \\ Y_n(t)=\stackrel{~}{Y}_n(t),t\underset{i=\mathrm{}+q+p+1}{\overset{n1}{}}\left(\text{supp}\mathrm{\Lambda }_i[0,a]\right).\hfill \end{array}$$

The claims of Step 4 are also easy to verify using Proposition 2.1, taking into account that we assume $`\mathrm{\Lambda }`$ is in basic form, which implies that $`\stackrel{~}{Y}_n(t^{})`$ cannot be a constant.

Another application of the ideas in the proof of Theorem 4.5 provides the following “Markov process” property.

###### Theorem 4.11

Fix $`0a1`$ such that $`a\text{supp}\mathrm{\Lambda }`$; the process $`Y(t)`$ satisfies

$$\sigma \{Y(t),t[0,a]\}\underset{\sigma \{Y(a)\}}{\text{ }\text{ }\text{ }\text{ }\text{ }}\sigma \{Y(t),t[a,1]\}$$

if and only if for all $`i`$, either $`\text{supp}\mathrm{\Lambda }_i[0,a[`$, or $`\text{supp}\mathrm{\Lambda }_i]a,1]`$.

Proof: One can use the same machinery as in the proof of Theorem 4.5: The ‘only if’ part can be proved within the same lines as Proposition 4.9, whereas for the other implication, if $`\text{supp}\mathrm{\Lambda }_i[0,a[`$, $`i=1,\mathrm{},\mathrm{}`$, and $`\text{supp}\mathrm{\Lambda }_i]a,1]`$, $`i=\mathrm{}+1,\mathrm{},n`$, then one can take

$$\begin{array}{c}Z^1:=(\stackrel{~}{Y}_1(a),\mathrm{},\stackrel{~}{Y}_n\mathrm{}(a))IR^n\mathrm{},\hfill \\ Z^2:=(\stackrel{~}{Y}_{n\mathrm{}+1}(a),\mathrm{},\stackrel{~}{Y}_n(a))IR^{\mathrm{}},\hfill \end{array}$$

and define $`g_1`$ as a function of $`Z^2`$ and the increments of the Wiener process in $`[0,a]`$, and $`g_2`$ as a function of $`Z^1`$ and the increments of the Wiener process in $`[a,1]`$.

###### Remark 4.12

Property (4.18) implies, using Lemma 3.3, that

$$\sigma \{Y(t),t[a,b]\}\underset{}{\text{ }\text{ }\text{ }\text{ }\text{ }}\sigma \{Y(t),t]a,b[^c\},$$

where $`=\sigma \{Y(a),Y(b);Y(t),t\text{supp}\mathrm{\Lambda }[a,b]\}`$. We conjecture that this property holds true for a linear functional boundary operator $`\mathrm{\Lambda }`$ supported on any subset of $`[0,1]`$, and that it is false in general if $``$ is replaced by a smaller $`\sigma `$-field. We will show a simple example illustrating the conjecture. Unfortunately, the technique we have employed here does not allow us to prove it.

Consider the process $`X(t):={\displaystyle _0^1}W(u)𝑑u+W(t)`$, solution of the first order problem

$$\{\begin{array}{c}\dot{X}(t)=\dot{W}(t),t[0,1]\hfill \\ _0^1X(u)𝑑u=0,\hfill \end{array}$$

in which the support of the boundary operator is the whole interval $`[0,1]`$. Fix $`a]0,1[`$, and set

$$𝒢:=\sigma \{X(u),u[0,a];X(1)\}\text{ and }:=\sigma \{X(u),u[a,1]\}.$$

We have trivially $`𝒢\underset{}{\text{ }\text{ }\text{ }\text{ }\text{ }}`$, but the conditioning $``$ cannot be replaced by the smaller $`\sigma `$-field $`^{}:=\sigma \{X(u),u[a,1][s,t]\}`$. Indeed, it is easy to see that

$$T:=\text{E}\left[_0^aX(u)𝑑u|\right]=_a^1X(u)𝑑u,$$

and one can check that $`T`$ is not $`^{}`$-measurable: Choose $`\omega ^1,\omega ^2C_0([0,1];IR)`$ such that $`\omega ^1\omega ^2`$ on $`]s,t[^c`$ and $`\omega ^1<\omega ^2`$ on $`]s,t[`$. An easy computation gives

$$T(\omega ^2)T(\omega ^1)=a_s^t(\omega ^2\omega ^1)(u)𝑑u>0.$$

Now, due to the continuity of $`T`$ as a functional on $`C_0([0,1];IR)`$, it is possible to find two open balls, centred at $`\omega ^1`$ and $`\omega ^2`$, such that their images through $`T`$ take values in two disjoint intervals of $`IR`$. We conclude that $`𝒢`$ and $``$ are not conditionally independent given $`^{}`$.


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# Event-by-event fluctuations in Mean 𝑝_𝑇 and Mean 𝑒_𝑇 in √𝑠_{𝑁𝑁}=130GeV Au+Au Collisions

## I Introduction

Phase instabilities near the QCD phase transition can result in non-statistical fluctuations that are detectable in final state observables . These instabilities, which may occur due to random color fluctuations , critical behavior at the QCD tri-critical point , or fluctuations from the decay of a Polyakov loop condensate , can result in a broadening of the transverse momentum or transverse energy distributions of produced particles for different classes of events. This phenomenon is expected to be detected experimentally by searching for deviations of the distributions of the event-by-event mean transverse momentum $`M_{p_T}`$ or mean transverse energy $`M_{e_T}`$ of produced particles from the random distributions expected for statistically independent particle emission.

An event-by-event analysis of $`M_{p_T}`$ was previously performed for $`158A`$GeV/c Pb+Pb Collisions at the CERN SPS by Experiment NA49 . In that analysis, the $`M_{p_T}`$ distributions measured over the rapidity range $`4<y_\pi <5.5`$ and $`p_T`$ range $`0.005<p_T<1.5`$GeV/c, were found to be consistent with random fluctuations. NA49 also performed an event-by-event analysis of the $`K/\pi `$ ratio , showing only very small deviations from random fluctuations. With an increase of $`\sqrt{s_{NN}}`$ to $`130`$GeV in RHIC collisions, unprecedented energy densities have been observed , hence conditions may be more favorable for a phase transition from hadronic matter to a Quark-Gluon Plasma which may be indicated in non-random fluctuations. Presented here is an event-by-event analysis of $`M_{p_T}`$ fluctuations and the first measurement of $`M_{e_T}`$ fluctuations at mid-rapidity at RHIC.

## II Analysis

The PHENIX experiment  consists of four spectrometers designed to measure simultaneously hadrons, leptons, and photons produced in nucleus-nucleus, proton-nucleus, and proton-proton collisions at RHIC. The two central arm spectrometers, which are located within a focusing magnetic field and each cover $`\pm 0.35`$ in pseudorapidity and $`\mathrm{\Delta }\varphi =90^{}`$ in azimuthal angle, are utilized in this analysis. The primary interaction trigger was defined using the Beam-Beam Counters (BBC)  and Zero Degree Calorimeters (ZDC) . Events are selected with a requirement that the collision vertex along the beam axis has $`|z|<20`$cm as measured by both the BBC and ZDC. Event centrality is defined using correlations in the BBC and ZDC analog response as described in . For the present analysis, the events are classified according to the $`05\%`$, $`010\%`$, $`1020\%`$, and $`2030\%`$ most central events.

The drift chamber  is used in conjunction with the innermost pad chamber, called PC1, to measure the transverse momentum of charged particles traversing the PHENIX acceptance. A fiducial section of the drift chamber is chosen to minimize the effect of time-dependent variations in the performance of the detector during the data-taking period. The fiducial volume of the $`M_{p_T}`$ analysis spans an azimuthal range of $`\mathrm{\Delta }\varphi =58.5^{}`$ and covers the pseudorapidity range $`|\eta |<0.35`$. Reconstructed tracks  are required to contain a match to a hit in PC1 to insure that the tracks are well reconstructed in three dimensions for reliable momentum determination.

The $`M_{e_T}`$ distribution is determined from clusters reconstructed in the two instrumented sectors of the Lead-Scintillator electromagnetic calorimeter . The quantity $`e_T`$ is defined as the transverse energy per reconstructed calorimeter cluster as described in , which can include clusters that have been merged. The effects of cluster merging on the $`M_{e_T}`$ distribution are discussed later. The fiducial volume of the $`M_{e_T}`$ analysis spans an azimuthal range of $`\mathrm{\Delta }\varphi =45^{}`$ and covers $`|\eta |<0.35`$.

There are no acceptance nor efficiency corrections applied to the semi-inclusive $`p_T`$ or $`e_T`$ distributions prior to the calculation of $`M_{p_T}`$ or $`M_{e_T}`$. These corrections do not vary from event to event and are identical for data and mixed events; therefore they do not modify the values of the fluctuation quantities defined later. The $`M_{p_T}`$ distributions are calculated using the formula

$$M_{p_T}=(1/N_{tracks})\underset{i=1}{\overset{N_{tracks}}{}}p_{T}^{}{}_{i}{}^{},$$
(1)
where $`N_{tracks}`$ is the number of tracks in the event that pass the cuts outlined above and lie within the $`p_T`$ range $`0.2<p_T<1.5`$GeV/c. Similarly, the $`M_{e_T}`$ distributions are calculated using the formula

$$M_{e_T}=(1/N_{clus})\underset{i=1}{\overset{N_{clus}}{}}e_{T}^{}{}_{i}{}^{},$$
(2)
where $`N_{clus}`$ is the number of calorimeter clusters in the event that lie within the $`e_T`$ range $`0.225<e_T<2.0`$GeV. An event is excluded from the analysis if $`N_{tracks}`$ or $`N_{clus}`$ is below a minimum value to insure that there are a sufficient number of tracks or clusters to determine a mean and to exclude background events. This minimum value for the $`05\%`$, $`010\%`$, $`1020\%`$, and $`2030\%`$ centrality classes, respectively, is 40, 30, 20, and 10 for the $`M_{p_T}`$ analysis and 30, 20, 10, and 10 for the $`M_{e_T}`$ analysis. Table I lists statistics pertaining to the data samples used to determine $`M_{p_T}`$ and Table II lists the statistics pertaining to the data samples used to determine $`M_{e_T}`$. The events used for the $`M_{p_T}`$ and $`M_{e_T}`$ analyses are considered independently of each other.

In order to compare the $`M_{p_T}`$ and $`M_{e_T}`$ distributions to what is expected for statistically independent particle emission, mixed events are considered as the baseline for the random distribution. To obtain a precision comparison, it is important to match the number of tracks or clusters along with the mean of the semi-inclusive distribution of the mixed events to that of the data. Therefore, in both analyses, mixed events are constructed by pre-determining the number of charged particle tracks or calorimeter clusters in the mixed event $`N_{mix}`$ by directly sampling the corresponding data $`N_{tracks}`$ or $`N_{clus}`$ distributions. Figure 1 shows a comparison of the $`N_{tracks}`$ distributions from the data and the normalized mixed event $`N_{mix}`$ distribution for the $`010\%`$ centrality class. Once $`N_{mix}`$ is determined, a mixed event is filled with $`p_T`$ or $`e_T`$ values from the data with the following criteria: a) no two $`p_T`$ or $`e_T`$ values from the same data event are allowed to reside in the same mixed event, b) only $`p_T`$ or $`e_T`$ values passing all cuts in the determination of $`M_{p_T}`$ or $`M_{e_T}`$ from the data events are placed in a mixed event, and c) only data events from the same centrality class are used to construct a mixed event corresponding to that class. Once a mixed event is filled with $`N_{mix}`$ tracks or clusters, its $`M_{p_T}`$ or $`M_{e_T}`$ is calculated in the same manner as for the data events.

For both analyses, the data contain a fraction of tracks or clusters within close physical proximity that have merged into a single track or cluster. This fraction is estimated by embedding simulated single-particle events that are processed through a detailed simulation of the detector response into real data events, which are then reconstructed in the same manner as the data. For the $`05\%`$ centrality class, we estimate that 6% of the tracks and 5% of the clusters are affected.

For the $`M_{p_T}`$ analysis, tracks that are merged into a single reconstructed track typically have similar values of $`p_T`$. The result is a slightly lower value of $`N_{tracks}`$ which causes a slight broadening in the width of the $`M_{p_T}`$ distribution due to the reduced statistics per event. However, since the $`N_{tracks}`$ data distribution is directly sampled during the construction of mixed events, the effect of merged tracks cancels for comparisons between the data and mixed events.

For the $`M_{e_T}`$ analysis, the effect of merged clusters is complicated by the fact that a single cluster is reconstructed with an $`e_T`$ corresponding to the sum of the two (or more) particles contributing to the cluster. To understand this effect on the mixed events, we note that the fraction of merged clusters within a data event increases with event multiplicity. Also, many of the data events with the lowest $`M_{e_T}`$ coincide with the lowest multiplicity events since they contain few, if any, merged clusters that would yield a higher $`M_{e_T}`$. When the merged clusters in the data events are randomly redistributed among the mixed events, low multiplicity mixed events can contain more merged clusters than the data events with the same multiplicity, resulting in a gross upward shift in $`M_{e_T}`$ for those mixed events. This results in apparent excess non-random fluctuations at low $`M_{e_T}`$. Conversely, high multiplicity mixed events can contain fewer merged clusters than the data events with the same multiplicity, resulting in a gross downward shift in $`M_{e_T}`$ for those mixed events. However, since the mean is taken over more clusters in this case, the effective shift in $`M_{e_T}`$ is reduced at high $`M_{e_T}`$, and the apparent non-random fluctuations are much less pronounced. An estimate of the magnitude of this effect is presented later.

## III Results

To compare directly the semi-inclusive $`p_T`$ distribution to the $`M_{p_T}`$ distribution assuming statistically independent particle emission, the closed form prescription outlined in  is used. This prescription describes the semi-inclusive $`p_T`$ distribution using a Gamma distribution,

$$f(p_T)=f_\mathrm{\Gamma }(p_T,p,b)=\frac{b}{\mathrm{\Gamma }(p)}(bp_T)^{p1}e^{bp_T},$$
(3)
where $`p`$ and $`b`$ are free parameters that are related to the mean and standard deviation of the semi-inclusive distribution as

$$p=\frac{<p_T>^2}{\sigma _{p_T}^2},b=\frac{<p_T>}{\sigma _{p_T}^2},$$
(4)
where

$$\sigma _{p_T}=(<p_T^2><p_T>^2)^{1/2}.$$
(5)
The reciprocal of $`b`$ is the inverse slope parameter of the $`p_T`$ distribution. With the track multiplicity distribution assumed to be a negative binomial distribution, $`f_{NBD}(N_{tracks},1/k,<N_{tracks}>)`$, the $`M_{p_T}`$ distribution can be calculated using

$$g(p_T)=\underset{N=N_{min}}{\overset{N_{max}}{}}f_{NBD}(N,1/k,<N>)f_\mathrm{\Gamma }(p_T,Np,Nb),$$
(6)
where the loop is over $`N_{tracks}`$ from $`N_{min}`$ to $`N_{max}`$, which are the limits of the multiplicity. The value of the negative binomial distribution parameter $`k`$ is given by

$$\frac{1}{k}=\frac{\sigma _{p_T}^2}{<N_{tracks}>^2}\frac{1}{<N_{tracks}>}.$$
(7)
Therefore, given $`<p_T>`$, $`\sigma _{p_T}`$, and $`<N_{tracks}>`$ extracted from the semi-inclusive $`p_T`$ distribution, the expected random $`M_{p_T}`$ distribution can be calculated. Figure 2 shows the $`M_{p_T}`$ distribution for the $`05\%`$ centrality class. Overlayed on the data as a dotted curve is the result of the calculation. The agreement between the data distribution and the calculation illustrates the absence of large non-statistical fluctuations in the data. The remainder of this article will quantify the amount of non-statistical fluctuations observed and place limits on the level of fluctuations that can be present in central Au+Au collisions at $`\sqrt{s_{NN}}=130`$GeV.

To quantify the magnitude of the deviation of fluctuations from the expectation of statistically independent particle emission, the magnitude of the fluctuation $`\omega _T`$ in the transverse quantity $`M_T`$, representing $`M_{p_T}`$ or $`M_{e_T}`$, is defined as

$$\omega _T=\frac{(<M_T^2><M_T>^2)^{1/2}}{<M_T>}=\frac{\sigma _{M_T}}{<M_T>}.$$
(8)
The value of $`\omega _T`$ is calculated independently for the data distribution and for the baseline, or mixed event, distribution. The difference in the fluctuation from a random baseline distribution is defined as

$$d=\omega _{(T,data)}\omega _{(T,baseline)}.$$
(9)
The sign of $`d`$ is positive if the data distribution contains a correlation, such as Bose-Einstein correlations , when compared to the baseline distribution. The fraction of fluctuations which deviate from the expectation of statistically independent particle emission is given by

$`F_T`$ $`={\displaystyle \frac{(\omega _{(T,data)}\omega _{(T,baseline)})}{\omega _{(T,baseline)}}}`$ (10)
$`={\displaystyle \frac{(\sigma _{(T,data)}\sigma _{(T,baseline)})}{\sigma _{(T,baseline)}}},`$ (11)
where $`\sigma _{(T,data)}`$ refers to the standard deviation of the event-by-event $`M_T`$ data distribution and $`\sigma _{(T,baseline)}`$ is the corresponding quantity for the baseline, or mixed event, distribution. In the absence of a common language for the analysis of $`M_{p_T}`$ and $`M_{e_T}`$ fluctuations, the commonly used fluctuation quantity $`\varphi _T`$  is also presented in order to compare this measurement to previous results . The quantity $`d`$ is related directly to $`\varphi _T`$ via

$`\varphi _T`$ $`=(\sigma _{(T,data)}\sigma _{(T,baseline)})\sqrt{<N_T>}`$ (12)
$`=d<M_T>\sqrt{<N_T>},`$ (13)
where $`N_T`$ represents $`N_{tracks}`$ or $`N_{clus}`$. The quantity $`\varphi _T`$ is related to $`F_T`$ by

$$\varphi _T=F_T\sigma _{(T,baseline)}\sqrt{<N_T>}.$$
(14)
The standard deviation of the semi-inclusive spectra can be approximated by $`\sigma _{(T,incl.)}\sigma _{(T,baseline)}\sqrt{<N_T>}`$ , where $`\sigma _{(T,incl.)}`$ is the standard deviation of the semi-inclusive distribution as defined in Eq. 3.3. Therefore, $`\varphi _T`$ is simply the fraction of non-random fluctuations in the event-by-event mean $`p_T`$ or $`e_T`$, $`F_T`$, scaled by $`\sigma _{(T,incl.)}`$. An advantage of $`F_T`$ over $`\varphi _T`$ is that measurements expressed in $`F_T`$ can be directly compared without further scaling.

The magnitudes of any non-random fluctuations are established by comparing the data distributions to the mixed event distributions, which serve as the random baseline distributions. For this purpose, the mixed event distributions are normalized to minimize the $`\chi ^2`$ value with respect to the data distributions. Figure 3 and Figure 5 show the $`M_{p_T}`$ and $`M_{e_T}`$ distributions for all four centrality classes (data points) with the corresponding mixed event $`M_{p_T}`$ and $`M_{e_T}`$ distributions overlayed on the data as dotted curves. The broadening of the distributions for less central collisions are due to the reduction in $`<N_{tracks}>`$ or $`<N_{clus}>`$. Shown in Figure 4 and Figure 6 are the residuals between the data and mixed events, defined for each bin $`i`$ as $`residual_i=(M_{(T,data)_i}M_{(T,mixed)_i})/\sigma _i`$, in units of standard deviations, for each centrality class. The shapes of the residual distributions are primarily driven by the normalization procedure applied to the mixed events.

For the $`M_{p_T}`$ distributions, the data and mixed event distributions are indistinguishable. However, the upper $`M_{e_T}`$ edges of the data and mixed event $`M_{e_T}`$ distributions show good agreement while the lower $`M_{e_T}`$ edge of the data distributions are slightly wider than the mixed event distribution. If this low $`e_T`$ effect were physical, it would imply fluctuations with slightly more low $`e_T`$ photons since the effect is not seen in the $`M_{p_T}`$ distribution for charged particle tracks. However, some of the excess fluctuations at low $`e_T`$ can be attributed to the effects of cluster merging previously discussed. The magnitude of this effect has been investigated using a Monte Carlo simulation which calculates $`M_{e_T}`$ after reproducing the calorimeter cluster separation distribution, the $`N_{clus}`$ distribution, and the semi-inclusive $`e_T`$ distributions from the data. The fluctuations in the $`M_{e_T}`$ distribution with this effect included in each event are compared to a simulated mixed event $`M_{e_T}`$ distribution constructed from the same generated dataset using the same procedure that is applied to the data. In this manner, it is estimated that the cluster merging effect contributes an additional $`F_T`$ = 1.5%, 2.1%, 0.9%, and less than 0.01% to the non-random fluctuations for the $`05\%`$, $`010\%`$, $`1020\%`$, and $`2030\%`$ centrality classes, respectively. The simulation confirms that the cluster merging effect significantly contributes only to the lower $`M_{e_T}`$ edge of the distribution. The remainder of the excess low $`e_T`$ fluctuations is likely due to correlated low energy background. GEANT  simulations indicate that the primary background contribution is produced by low energy electrons and muons that scatter off of the pole tips of the central arm spectrometer magnet but still pass the cluster selection cuts. Because of the difficulty in quantifying the contribution of background to the excess fluctuations, the present $`M_{e_T}`$ data are taken to indicate an upper limit on non-statistical fluctuations rather than an indication of true non-statistical fluctuations.

The values of $`\omega _T`$, $`d`$, $`F_T`$, and $`\varphi _T`$ for each centrality class using the mixed events as the random baseline distribution are tabulated in Table III for $`M_{p_T}`$ and Table IV for $`M_{e_T}`$. The errors quoted for these quantities include statistical errors and systematic errors due to time-dependent variations over the data-taking period. The systematic errors are estimated by dividing each dataset into nine subsets with each subset containing roughly equal numbers of events. For the $`M_{p_T}`$ analysis, the systematic errors contribute to 81%, 88%, 76%, and 75% of the total error in $`\omega _T`$ and 85%, 88%, 80%, and 85% of the total error in the variables $`d`$, $`F_T`$, and $`\varphi _t`$ for the $`05\%`$, $`010\%`$, $`1020\%`$, and $`2030\%`$ centrality classes, respectively. The corresponding values for the $`M_{e_T}`$ analysis are a 67%, 63%, 81%, and 82% contribution to the total errors in $`\omega _T`$, and a 64%, 63%, 81%, and 82% contribution to the total errors in $`d`$, $`F_T`$, and $`\varphi _t`$ for each centrality class. The cluster merging contribution estimates noted above are not applied to the values quoted in Table IV.

## IV Discussion

Based upon the fluctuation measurements presented here, certain fluctuation scenarios in RHIC Au+Au collisions at $`\sqrt{s_{NN}}=130`$GeV are excluded. For this purpose, we consider two variations of a model that contains two classes of events with a difference of effective temperature, defined as $`\mathrm{\Delta }T=T_2T_1`$, where $`T_2`$ is the inverse slope parameter of the event class with the higher effective temperature, and $`T_1`$ is the inverse slope parameter of the event class with the lower effective temperature. The first variation, Model A, will consider a case where the means of the semi-inclusive $`p_T`$ spectra for the two event classes are identical, but the standard deviations are different. The second variation, Model B, will consider a case where the means of the semi-inclusive $`p_T`$ spectra are different, but the standard deviations are identical. Since the semi-inclusive $`p_T`$ distribution is an observed quantity, the two event classes must be constrained in such a way that the mean and standard deviation of the final semi-inclusive $`p_T`$ distribution remains constant while the effect of the fluctuation manifests itself in the $`M_{p_T}`$ distribution.

The dual event class model is applied to the determination of the sensitivity to fluctuations in $`M_{p_T}`$ for the $`05\%`$ centrality class as follows. Returning to the prescription outlined in , the semi-inclusive transverse $`p_T`$ spectrum can be parameterized by the $`f_\mathrm{\Gamma }(p_T,p,b)`$ distribution defined in Eq. 3.1. For both model variations, the fraction of events in the event class with the higher effective temperature is defined as

$$q=\frac{(N_{events})_{class1}}{(N_{events})_{class1}+(N_{events})_{class2}}.$$
(15)
The $`p_T`$ distribution of the combined sample can then be expressed as

$$f(p_T)=q\mathrm{\Gamma }(p_T,p_1,b_1)+(1q)\mathrm{\Gamma }(p_T,p_2,b_2),$$
(16)
where $`T_1=1/b_1`$ and $`T_2=1/b_2`$.

For Model A, the semi-inclusive $`p_T`$ distributions of each event class are constrained to have the same mean, so we require

$$\mu =p/b=p_1/b_1=p_2/b_2.$$
(17)
The variance of the final semi-inclusive $`p_T`$ distribution for Model A is constrained by

$$\frac{\sigma ^2}{\mu ^2}=\frac{1}{p}=\frac{q}{p_1}+\frac{(1q)}{p_2}.$$
(18)
With these constraints, the choice of a value for $`q`$ and the effective temperature of one event class is sufficient to extract the remaining parameters from which sensitivity estimates for fluctuations in $`M_{p_T}`$ are obtained.

For Model B, the semi-inclusive $`p_T`$ distributions of each event class are allowed to have different means, $`\mu _1`$ and $`\mu _2`$, so the mean of the total semi-inclusive distribution can be expressed as $`\mu =q\mu _1+(1q)\mu _2`$. Defining a mean shift, $`\mathrm{\Delta }\mu `$, as $`\mathrm{\Delta }\mu =\mu _2\mu _1`$, we obtain

$$\mu _2=\mu +q\mathrm{\Delta }\mu .$$
(19)
Allowing $`p_1`$ = $`p_2`$ and applying the constraint that the variances of the two event classes are identical, yields

$$\frac{1}{p_1}=\frac{\frac{1}{p}q(1q)(\frac{\mathrm{\Delta }\mu }{\mu })^2}{1+q(1q)(\frac{\mathrm{\Delta }\mu }{\mu })^2}.$$
(20)
With a choice of values for $`q`$ and $`\mathrm{\Delta }\mu `$, the remaining parameters can be calculated, including $`\mathrm{\Delta }T`$.

Both variations of the dual event class model are implemented in a Monte Carlo simulation in the following manner. The number of particles in an event is determined by sampling the $`N_{tracks}`$ data distribution, approximated by a Gaussian distribution fit to the data. The $`p_T`$ of each particle in an event is determined individually by sampling the appropriate $`\mathrm{\Gamma }(p_T,p,b)`$ distribution fit to the semi-inclusive $`p_T`$ data distribution, which yields $`p=0.8`$ and $`b=2.46`$ for $`05\%`$ centrality. The $`p_T`$ of each particle is restricted to the $`p_T`$ range of the measurement. With $`N_{tracks}`$ and the $`p_T`$ distribution determined, the $`M_{p_T}`$ for a given number of events is calculated. The generated $`M_{p_T}`$ distribution with $`q=0`$ for either model variation is found to be statistically consistent with the mixed event $`M_{p_T}`$ distribution.

The data contain a fraction of background particles that did not originate from the collision vertex that effectively dilute the sensitivity to non-random fluctuations. To address this, a fraction of the particles in an event are randomly tagged as background particles, whose $`p_T`$ distribution is then generated with a separate parameterization prior to calculating $`M_{p_T}`$ for an event. The level of background contamination is estimated by processing HIJING  Au+Au events through a software chain that includes a detailed GEANT simulation  with the complete PHENIX detector geometry included, followed by a detailed simulation of the detector electronics response , whose output is then processed by track, cluster, and momentum reconstruction using the identical software and input parameters as is used for the data analysis. It is estimated that 11% of the tracks and 26% of the clusters are due to background particles, independent of centrality class over the centrality range of these measurements. The estimated $`p_T`$ and $`e_T`$ distributions for the background particles are well parameterized by exponential distributions. Again, the majority of the $`e_T`$ background occurs at low $`e_T`$, so any correlated background would most likely contribute to the lower side of the $`M_{e_T}`$ distribution.

To determine the sensitivity to fluctuations within the dual event class model, the fluctuation fraction, $`q`$, and the value of $`p_1`$ for Model A and $`\mathrm{\Delta }\mu `$ for Model B are varied and the $`M_{p_T}`$ distribution is generated at each step. A chi-square test is then performed on the generated $`M_{p_T}`$ distribution with respect to the mixed event data $`M_{p_T}`$ distribution. For a given value of $`q`$, the $`\chi ^2`$ result increases as $`\mathrm{\Delta }T`$ increases, which allows a fluctuation exclusion region to be defined for the single degree of freedom. The curves in Figure 7 show the lower exclusion boundaries for the $`05\%`$ centrality $`M_{p_T}`$ measurement at the 95% Confidence Level as a function of $`q`$ and $`\mathrm{\Delta }T`$ for both variations of the model. If the sensitivity is determined based upon the non-mixed data distribution, the lower exclusion boundary increases by less than 2 MeV for all values of $`q`$ for either model. Also, for all values of $`q`$ in either model, the estimated background contribution degrades the sensitivity estimates by $`\mathrm{\Delta }T`$ = 3 MeV for both models.

A recent model of event-by-event fluctuations where the temperature parameter $`T=1/b`$ fluctuates with a standard deviation $`\sigma _T`$ on an event-by-event basis , can be simply related to $`F_T`$:

$$\frac{\sigma _T^2}{<T>^2}=\frac{2F_T}{p(<n>1)},$$
(21)
where $`p=0.8`$ is the semi-inclusive parameter extracted from the present data. For the $`05\%`$ centrality class, the present measurement establishes a 95% Confidence Limit of 2.6 $`\times 10^3`$ for $`\sigma _T^2/<T>^2`$, or 5% for $`\sigma _T/<T>`$.

## V Conclusions

The fluctuations in the $`M_{p_T}`$ distributions for all centrality classes are consistent with the presence of no fluctuations in excess of the random expectation. The magnitude of $`F_T`$ in all cases is positive, which may be due to the presence of Hanbury-Brown-Twiss correlations. The fluctuations in the $`M_{e_T}`$ distributions do have a small non-statistical component, much of which is attributable to the effects of merged clusters, the remainder of which are taken to indicate an upper limit on non-statistical fluctuations in transverse energy. By defining a dual event class model, limits are set on the amount of $`M_{p_T}`$ fluctuations that can be present in the angular aperture $`|\eta |<0.35`$ and $`\mathrm{\Delta }\varphi =58.5^{}`$ in $`\sqrt{s_{NN}}=130`$GeV Au+Au collisions. During the RHIC run of 2001, PHENIX has taken data for $`\sqrt{s_{NN}}=200`$GeV Au+Au collisions with about a factor of four increase in azimuthal angular acceptance for both the $`M_{p_T}`$ and $`M_{e_T}`$ analyses, which will allow the measurements to be extended toward more peripheral collisions.

## VI Acknowledgements

We thank the staff of the RHIC Project, Collider-Accelerator, and Physics Departments at Brookhaven National Laboratory and the staff of the other PHENIX participating institutions for their vital contributions. We acknowledge support from the Department of Energy, National Science Foundation, and Dean of the College of Arts and Sciences, Vanderbilt University (U.S.A), Ministry of Education, Culture, Sports, Science, and Technology and the Japan Society for the Promotion of Science (Japan), Russian Academy of Science, Ministry of Atomic Energy of Russian Federation, Ministry of Industry, Science, and Technologies of Russian Federation (Russia), Bundesministerium fuer Bildung und Forschung, Deutscher Akademischer Auslandsdienst, and Alexander von Humboldt Stiftung (Germany), VR and the Wallenberg Foundation (Sweden), MIST and the Natural Sciences and Engineering Research Council (Canada), Conselho Nacional de Desenvolvimento Científico e Tecnológico and Fundação de Amparo à Pesquisa do Estado de São Paulo (Brazil), Natural Science Foundation of China (People’s Republic of China), IN2P3/CNRS (France), Department of Atomic Energy and Department of Science and Technology (India), Korea Research Foundation and Center for High Energy Physics (Korea), the U.S. Civilian Research and Development Foundation for the Independent States of the Former Soviet Union, and the US-Israel Binational Science Foundation.