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# Weisskopf-Wigner Decay Theory for the Energy-Driven Stochastic Schrödinger Equation

## I Introduction

There has recently been considerable interest in the possibility that quantum mechanics, and the Schrödinger equation, may be modified at a very low level by effects arising from Planck scale physics. Such speculations have been motivated on the one hand by considerations from string theory ellis  and quantum gravity penrose , and on the other hand by efforts penrose ; pearle ; gisin  to achieve an objective equation describing state vector reduction. The majority of the objective reduction discussions fall into two classes: those that postulate a stochastic process producing spatial localization pearle , and those that postulate an analogous stochastic process leading to localization in energy gisin  (the so-called “energy-driven” stochastic Schrödinger equation.) Both the spatial localization and the energy localization stochastic Schrödinger equations avoid problems with superluminal signal propagation that characterize attempts at deterministic nonlinear modifications of the Schrödinger equation bialy . We find the energy-driven approach particularly appealing because it is energy conserving, leads with no approximations to Born rule probabilities and to the Lüders projection postulate, has sensible clustering properties, and when environmental interactions are taken into account explains state vector reduction with a single Planck scale stochastic parameter gisin ; adler .

Although physical prejudices might suggest a Planck scale magnitude for the stochastic parameter in the energy-driven equation, one can instead take the point of view that the stochastic parameter can have a priori any value, and use current experimental information to place bounds on it. This approach has been pursued lisi  in the context of particle physics systems that exhibit oscillations between different mass eigenstates (the $`K`$-meson, $`B`$-meson, and neutrino systems), with results that are summarized in the final section of this paper. An alternative source of bounds on the stochastic parameter could come from experiments observing decays and line shapes in atomic and particle systems, since if the stochastic terms in the Schrödinger equation were to change the standard Weisskopf-Wigner analysis of decay processes in a significant way, then observable effects could result. Thus, to pursue phenomenological studies of the energy-driven equation, it is important to generalize the standard Weisskopf-Wigner decay theory weiss  to include effects of the energy-driven stochastic terms. This is the problem that is analyzed in this paper.

## II The Energy-Driven Stochastic Schrödinger Equation and Properties of the Itô Stochastic Calculus

Letting $`|\psi `$ denote a unit normalized Schrödinger picture state vector, the standard form pearle ; gisin ; adler  of the energy-driven stochastic Schrödinger equation is (with $`\mathrm{}=1`$)

$$d|\psi =iH|\psi dt\frac{1}{8}\sigma ^2(HH)^2|\psi dt+\frac{1}{2}\sigma (HH)|\psi dW_t.$$
(1a)
Here $`H`$ is the Hamiltonian, $`H=\psi |H|\psi `$ is the expectation of the Hamiltonian in the state $`|\psi `$, $`\sigma `$ is a numerical parameter governing the strength of the stochasticity, and $`dW_t`$ is an Itô stochastic differential that, together with $`dt`$, obeys the standard Itô calculus rules gardiner
$$dW_t^2=dt,dW_tdt=dt^2=0.$$
(1b)
By construction, the nonlinear evolution of Eq. (1a) guarantees the preservation in time of the unit normalization of the state vector $`|\psi `$.

In the following sections, we shall need a number of properties of the Itô calculus that we summarize here. First of all, in the Itô calculus the Leibnitz chain rule generalizes to

$$d(AB)=(A+dA)(B+dB)AB=(dA)B+AdB+dAdB,$$
(2a)
with the final term $`dAdB`$ contributing a term proportional to $`dt`$ when the $`dW_t`$ contributions to both $`dA`$ and $`dB`$ are nonzero. Letting $`W_t`$ be the Brownian motion
$$W_t=_0^t𝑑W_u,$$
(2b)
we see in particular that
$$d\mathrm{exp}(\alpha W_t)=\mathrm{exp}(\alpha W_t)[\mathrm{exp}(\alpha dW_t)1]=\mathrm{exp}(\alpha W_t)[\alpha dW_t+\frac{1}{2}\alpha ^2dt].$$
(2c)
Letting $`E[\mathrm{}]`$ denote the stochastic expectation of its argument, and letting $`A(t)`$ denote any function of the stochastic process up to time $`t`$, we have

$$E[dW_tA(t)]=0,$$
(3a)
since the Itô differential refers to the time interval from $`t`$ to $`t+dt`$, and hence is statistically independent of the process up to time $`t`$. Thus, taking the expectation of Eq. (2c), we get the differential equation
$$dE[\mathrm{exp}(\alpha W_t)]=E[\mathrm{exp}(\alpha W_t)]\frac{1}{2}\alpha ^2dt,$$
(3b)
which can be immediately integrated to give
$$E[\mathrm{exp}(\alpha W_t)]=\mathrm{exp}(\frac{1}{2}\alpha ^2t),$$
(3c)
a result that will be needed later on.

Let us make an elementary application of the Itô formalism, to write the stochastic Schrödinger equation of Eq. (1a) in an equivalent form. First of all, forming the density matrix

$$\rho =|\psi \psi |,$$
(4a)
we have from Eq. (2a),
$$d\rho =(d|\psi )\psi |+|\psi d\psi |+d|\psi d\psi |,$$
(4b)
which on substitution of Eq. (1a) and use of the Itô calculus rules of Eq. (1b) gives the evolution equation for the density matrix,
$$d\rho =i[\rho ,H]dt\frac{1}{8}\sigma ^2[H,[H,\rho ]]dt+\frac{1}{2}\sigma [\rho ,[\rho ,H]]dW_t.$$
(4c)
Taking the stochastic expectation of this equation, using Eq. (3a), gives a differential equation of the Lindblad type lind  for $`E[\rho ]`$,

$$\frac{dE[\rho ]}{dt}=i[E[\rho ],H]\frac{1}{8}\sigma ^2[H,[H,E[\rho ]]].$$
(5)
The fact that this equation is linear (in contrast to Eq. (4c), which is nonlinear) is the fundamental reason bialy  why Eq. (1a) does not give rise to superluminal signal propagation.

## III Initial Formulation of the Decay Problem

Let us now formulate the decay problem for the stochastic Schrödinger equation of Eq. (1a), following the standard procedure for the usual Schrödinger equation without stochastic terms. We suppose that for times $`t0`$ the Hamiltonian $`H`$ is given by an unperturbed Hamiltonian $`H_0`$, with eigenstates $`|n`$ and eigenvalues $`E_n`$,

$$H_0|n=E_n|n,$$
(6a)
and that the system under consideration is in an eigenstate $`|s_A`$ with eigenvalue $`E_s`$, which is one of a set of degenerate energy eigenstates $`|s_a,a=1,\mathrm{},D`$. Because Eq. (1a) acts as an ordinary Schrödinger evolution on a state $`|\psi `$ that is an energy eigenstate, the system remains in the state $`|s_A`$ as long as the Hamiltonian remains equal to $`H_0`$. Hence the starting point for the standard decay analysis merz  is also a consistent starting point for its stochastic extension under Eq. (1a). As in the standard procedure, we assume that at $`t=0`$ a time-independent perturbation $`V`$ is switched on, so that for times $`t>0`$ the Hamiltonian is $`H=H_0+V`$. The initial state $`|s_A`$ is then no longer an energy eigenstate, and so will decay into various other states $`|m`$; our problem, as in the usual case, is to find the partial transition rates for this decay and the probability amplitude for the system to remain in the initial degenerate group of states.
In formulating this problem, it is convenient to expand the state $`|\psi `$ over the basis $`|n`$ and, as in the standard case, to remove the Schrödinger time evolution associated with the unperturbed Hamiltonian $`H_0`$, by writing

$$|\psi (t)=\underset{n}{}|n\mathrm{exp}(iE_nt)C_n(t).$$
(6b)
Substituting Eq. (6b) into Eq. (1a), and projecting on $`m|`$, it is a matter of straightforward but somewhat tedious algebra to compute the stochastic evolution equation for $`C_m(t)`$, with the result

$`dC_m(t)`$ $`=`$ $`\alpha _mC_m(t)+{\displaystyle \underset{n}{}}\beta _{mn}C_n(t),`$
$`\alpha _m`$ $`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(E_mH)^2dt+{\displaystyle \frac{1}{2}}\sigma (E_mH)dW_t,`$
$`\beta _{mn}`$ $`=`$ $`iV_{mn}\mathrm{exp}[i(E_mE_n)t]dt{\displaystyle \frac{1}{8}}\sigma ^2[(E_m+E_n2H)V_{mn}+(V^2)_{mn}]\mathrm{exp}[i(E_mE_n)t]dt`$ (7a)
$`+`$ $`{\displaystyle \frac{1}{2}}\sigma V_{mn}\mathrm{exp}[i(E_mE_n)t]dW_t.`$
The corresponding expression for $`H`$ is
$$H=\underset{n}{}E_n|C_n(t)|^2+\underset{mn}{}V_{mn}\mathrm{exp}[i(E_mE_n)t]C_m^{}(t)C_n(t).$$
(7b)
In these equations $`V_{mn}`$ and $`(V^2)_{mn}`$ denote the respective matrix elements
$$V_{mn}=m|V|n,(V^2)_{mn}=m|V^2|n.$$
(7c)
## IV Approximation to Leading Order in $`V`$

Equations (7a-c) are a complicated, nonlinear set of stochastic differential equations, and so to solve them approximations will be needed. Following the Weisskopf-Wigner analysis, we shall make the approximation of regarding $`V`$ as a small perturbation. The coefficients $`C_m,m\{s_a\}`$ for states not in the initial degenerate manifold will then be of order $`O(V)`$, and we neglect $`O(V^2)`$ and higher contributions to them (except those arising implicitly through our solution for the $`C_{s_a}`$). On the other hand, the coefficients $`C_{s_a}`$ of states in the degenerate manifold can be of order unity, and we calculate these coefficients to order $`V^2`$ accuracy, neglecting corrections of order $`V^3`$ and higher. In a similar fashion, in expressions involving the stochasticity parameter $`\sigma `$, we shall retain terms of order $`\sigma V`$ and its powers $`(\sigma V)^2`$, etc., but shall neglect terms of order $`\sigma V^2`$ and higher that involve extra factors of $`V`$ relative to the terms that we are retaining. Finally, although we shall see that $`E_mE_s`$ is effectively small, we shall retain all terms of order $`\sigma (E_mE_s)`$, $`\sigma ^2(E_mE_s)`$, $`[\sigma (E_mE_s)]^2`$, etc., but shall drop terms $`\sigma (E_mE_s)\sigma O(V^2)`$ that are smaller than these by a factor of order $`\sigma V^2`$ or $`V^2`$.

Making use of the perturbative ordering of the coefficients $`C_n`$, we begin by simplifying and approximating the expression in Eq. (7b) for $`H`$. Separating off the states in the initial degenerate manifold, the sum in Eq. (7b) becomes

$$H=E_s\underset{a}{}|C_{s_a}|^2+\underset{ab}{}V_{s_as_b}C_{s_a}^{}C_{s_b}+O(V^2).$$
(8a)
However, since the state vector $`|\psi `$ remains unit normalized, we have
$$\underset{a}{}|C_{s_a}|^2=1\underset{m\{s_a\}}{}|C_m|^2=1+O(V^2),$$
(8b)
and so we have
$$H=E_s+\underset{ab}{}V_{s_as_b}C_{s_a}^{}C_{s_b}+O(V^2).$$
(8c)
If we substitute Eq. (8c) back into Eq. (7a), we are still left with a nonlinear set of equations. Therefore we shall also introduce the simplifying assumption that the perturbing potential has vanishing matrix elements within the degenerate manifold containing the initial state, so that

$$V_{s_as_b}=0,a,b=1,\mathrm{},D.$$
(9a)
There are important, physically relevant cases in which Eq. (9a) is obeyed as a result of selection rules. For example, for radiative decays treated in the electric dipole approximation, with $`H_0`$ taken as the atomic Hamiltonian plus the free radiation Hamiltonian, and with $`V`$ taken as the atomic coupling to the transverse electromagnetic modes, Eq. (9a) is obeyed as a result of parity invariance when the states in the initial degenerate manifold all have the same parity. (We caution, however, that Eq. (9a) is not valid for the analysis of $`K`$ (or $`B`$) meson systems when $`H_0`$ is taken as the strong interaction Hamiltonian. Thus here either one has to employ the nonlinear equations following from Eq. (8c), or one has to redefine $`H_0`$ so as to impose Eq. (9a) by including in $`H_0`$ the $`|\mathrm{\Delta }S|=2`$ (or $`|\mathrm{\Delta }C|=2`$) weak interaction effective Hamiltonian terms, with $`V`$ defined to contain only the $`|\mathrm{\Delta }S|=1`$ (or $`|\mathrm{\Delta }C|=1`$) weak interaction terms responsible for $`K`$ (or $`B`$) meson decays. Such a redefinition is consistent in the vacuum saturation approximation for the $`|\mathrm{\Delta }S|=2`$ (or $`|\mathrm{\Delta }C|=2`$) terms.)
With the simplifying assumption of Eq. (9a), Eq. (8c) becomes simply

$$H=E_s+O(V^2).$$
(9b)
Substituting this into Eq. (7a), and dropping terms that are not of leading order in $`V`$ in the sense defined above, Eqs. (7a,7b) simplify to the following set of linear equations,

$`dC_m(t)`$ $`=`$ $`(\alpha _m^{(1)}dW_t+\alpha _m^{(2)}dt)C_m(t)+{\displaystyle \underset{n}{}}\mathrm{exp}[i(E_mE_n)t](\gamma _{mn}^{(1)}dW_t+\gamma _{mn}^{(2)}dt)C_n(t),`$
$`\alpha _m^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma (E_mE_s),\alpha _m^{(2)}={\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2={\displaystyle \frac{1}{2}}(\alpha _m^{(1)})^2,`$
$`\gamma _{mn}^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma V_{mn},\gamma _{mn}^{(2)}=iV_{mn}{\displaystyle \frac{1}{8}}\sigma ^2[(E_m+E_n2E_s)V_{mn}+(V^2)_{mn}].`$ (10a)
Corresponding to the magnitude ordering of the coefficients $`C_m`$ introduced above, it is convenient to rewrite Eq. (10a) as separate equations for the two cases, $`m\{s_a\}`$ and $`m\{s_a\}`$. For $`m\{s_a\}`$ the coefficients $`\alpha _s^{(1,2)}`$ vanish; separating the sum over $`n`$ into terms where $`n\{s_a\}`$ and $`n\{s_a\}`$, using the assumption of Eq. (9a), and dropping terms of nonleading order in $`V`$, we get
$`dC_{s_a}(t)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2dt{\displaystyle \underset{b}{}}(V^2)_{s_as_b}C_{s_b}(t)`$
$`+`$ $`{\displaystyle \underset{n\{s_a\}}{}}\mathrm{exp}[i(E_sE_n)t](\gamma _{s_an}^{(1)}dW_t+\gamma _{s_an}^{(2)}dt)C_n(t),`$
$`\gamma _{s_an}^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma V_{s_an},\gamma _{s_an}^{(2)}iV_{s_an}f_n,`$ (10b)
where we have introduced the definition
$$f_n=1\frac{i}{8}\sigma ^2(E_nE_s).$$
(10c)
For $`m\{s_a\}`$ the coefficients $`\alpha _m^{(1,2)}`$ are nonzero, but only the terms with $`n\{s_a\}`$ have to be retained in the sum over $`n`$, and so we similarly get
$`dC_m(t)`$ $`=`$ $`(\alpha _m^{(1)}dW_t+\alpha _m^{(2)}dt)C_m(t)`$
$`+`$ $`\mathrm{exp}[i(E_mE_s)t]{\displaystyle \underset{a}{}}(\gamma _{ms_a}^{(1)}dW_t+\gamma _{ms_a}^{(2)}dt)C_{s_a}(t),`$
$`\gamma _{ms_a}^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma V_{ms_a},\gamma _{ms_a}^{(2)}iV_{ms_a}f_m.`$ (10d)
Equations (10a-d) are the basic system of stochastic differential equations that we shall solve in the subsequent sections.

## V Equations for Expectations of the Coefficients

The principal quantities that we wish to calculate are the expectations $`E[|C_m(t)|^2]`$ of the squared magnitudes of the coefficients, since these give the expectations of the probabilities for the various states to be occupied. We shall show in this section that, within our approximations, these can be directly related to the expectations $`E[C_m(t)]`$ of the coefficients themselves, for which we shall derive a closed, linear set of ordinary differential equations.

Again, we consider separately the cases $`m\{s_a\}`$ and $`m\{s_a\}`$. For $`C_{s_a}`$, we write

$$C_{s_a}(t)=E[C_{s_a}(t)]+\mathrm{\Delta }_a(t),$$
(11a)
with $`E[\mathrm{\Delta }_a(t)]=0`$, and with $`\mathrm{\Delta }_a(0)=0`$ since the stochastic terms in the differential equation act only after $`t=0`$. However, referring to Eq. (10b) we see that $`dC_{s_a}`$ is of order $`V^2`$, and so $`\mathrm{\Delta }_a(t)`$ must also be of order $`V^2`$. Therefore
$$E[|C_{s_a}(t)|^2]=|E[C_{s_a}(t)]|^2+O(V^4),$$
(11b)
and so to the accuracy to which we are working, we can compute $`E[|C_{s_a}(t)|^2]`$ from the expectation $`E[C_{s_a}(t)]`$, ignoring the effects of fluctuations.

We consider next $`E[|C_m(t)|^2]`$ for $`m\{s_a\}`$. Applying the Itô rule of Eq. (2a), we have

$$dE[|C_m(t)|^2]=E[(dC_m^{}(t))C_m(t)+C_m^{}(t)dC_m(t)+dC_m^{}(t)dC_m(t)].$$
(12a)
Substituting Eq. (10d) for $`dC_m(t)`$ and using Eq. (3a), which eliminates the $`dW_t`$ contributions, and using the fact that to leading order in $`V`$ we can replace $`C_{s_a}(t)`$ by its expectation, we get after some algebraic simplification the formula
$`{\displaystyle \frac{d}{dt}}E[|C_m(t)|^2]`$ $`=`$ $`\mathrm{exp}[i(E_mE_s)t]if_mE[C_m(t)]{\displaystyle \underset{a}{}}V_{ms_a}^{}E[C_{s_a}^{}(t)]`$ (12b)
$``$ $`\mathrm{exp}[i(E_mE_s)t]if_m^{}E[C_m^{}(t)]{\displaystyle \underset{a}{}}V_{ms_a}E[C_{s_a}(t)]`$
$`+`$ $`{\displaystyle \frac{1}{4}}\sigma ^2|{\displaystyle \underset{a}{}}V_{ms_a}E[C_{s_a}(t)]|^2,`$
which can be integrated to give $`E[|C_m(t)|^2]`$ once the expectations $`E[C_m(t)]`$ and $`E[C_{s_a}(t)]`$ are known.

To get a closed set of equations for the expectations of the coefficients, we simply take the expectations of Eqs. (10b) and (10d), and use Eq. (3a), which again eliminates the $`dW_t`$ contributions. For $`E[C_{s_a}]`$ we thus get

$`{\displaystyle \frac{d}{dt}}E[C_{s_a}(t)]`$ $`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2{\displaystyle \underset{b}{}}(V^2)_{s_as_b}E[C_{s_b}(t)]`$ (13a)
$`+`$ $`{\displaystyle \underset{n\{s_a\}}{}}\mathrm{exp}[i(E_sE_n)t](i)V_{s_an}f_nE[C_n(t)],`$
while for $`E[C_m(t)]`$ with $`m\{s_a\}`$ we find
$`{\displaystyle \frac{d}{dt}}E[C_m(t)]`$ $`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2)E[C_m(t)]`$ (13b)
$`+`$ $`\mathrm{exp}[i(E_mE_s)t]{\displaystyle \underset{b}{}}(i)V_{ms_b}f_mE[C_{s_b}(t)].`$
## VI Solutions for Expectations of the Coefficients

We proceed now to solve the linear system of equations for the expectations of the coefficients given in Eqs. (13a,b). Since the problem is defined on the half line $`t>0`$, the natural way to do this is by using the Laplace transform. Defining

$$g_m(p)=_0^{\mathrm{}}𝑑t\mathrm{exp}(pt)E[C_m(t)],$$
(14a)
we have, by an integration by parts,
$$_0^{\mathrm{}}𝑑t\mathrm{exp}(pt)\frac{dE[C_m(t)]}{dt}=pg_m(p)E[C_m(0)],$$
(14b)
and also
$$_0^{\mathrm{}}𝑑t\mathrm{exp}(pt)\mathrm{exp}[i(E_mE_n)t]E[C_n(t)]=g_n(piE_m+iE_n),$$
(14c)
with the integrals in Eqs. (14a-c) defining analytic functions of $`p`$ in the right hand half plane $`\mathrm{Re}p>0`$. The inversion of the Laplace transform is given by the formula
$$E[C_m(t)]=\frac{1}{2\pi i}_{ϵi\mathrm{}}^{ϵ+i\mathrm{}}𝑑p\mathrm{exp}(pt)g_m(p),$$
(14d)
with $`ϵ>0`$ an infinitesimal positive constant.

Taking the Laplace transform of Eqs. (13a,b), and using the initial conditions $`E[C_{s_a}(0)]=C_{s_a}(0)=\delta _{aA}`$ and $`E[C_m(0)]=C_m(0)=0,m\{s_a\}`$, we get

$`pg_{s_a}(p)\delta _{aA}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2{\displaystyle \underset{b}{}}(V^2)_{s_as_b}g_{s_b}(p)`$ (15a)
$`+`$ $`{\displaystyle \underset{n\{s_a\}}{}}(i)V_{s_an}f_ng_n(p+iE_niE_s),`$
and for $`m\{s_a\}`$,
$`pg_m(p)`$ $`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2)g_m(p)`$ (15b)
$`+`$ $`{\displaystyle \underset{b}{}}(i)V_{ms_b}f_mg_{s_b}(p+iE_siE_m).`$
Solving Eq. (15b) for $`g_m(p)`$, $`m\{s_a\}`$, and shifting $`pp+iE_m`$ in the solution, we get

$$g_m(p+iE_m)=[p+iE_m+\frac{1}{8}\sigma ^2(E_mE_s)^2]^1\underset{b}{}(i)V_{ms_b}f_mg_{s_b}(p+iE_s).$$
(16a)
Shifting $`pp+iE_s`$ in Eq. (15a), and then substituting Eq. (16a), we get an algebraic equation for the set of quantities $`g_{s_b}(p+iE_s)`$,
$`{\displaystyle \underset{b}{}}K_{ab}g_{s_b}(p+iE_s)`$ $`=`$ $`\delta _{aA},`$
$`K_{ab}=(p+iE_s)\delta _{ab}`$ $`+`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(V^2)_{s_as_b}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{f_m^2V_{s_am}V_{ms_b}}{p+iE_s+i(E_mE_s)f_m}}.`$ (16b)
In physically interpreting these equations, we must remember that the Laplace transform variable $`p`$ is related to the usual energy variable $`E`$ by $`p=iE`$. Making this substitution in Eqs. (16a,b) we have respectively

$$g_m(iE+iE_m)=[iE+iE_m+\frac{1}{8}\sigma ^2(E_mE_s)^2]^1\underset{b}{}(i)V_{ms_b}f_mg_{s_b}(iE+iE_s),$$
(17a)
and
$`{\displaystyle \underset{b}{}}K_{ab}g_{s_b}(iE+iE_s)=\delta _{aA},`$
$`K_{ab}=(iE+iE_s)\delta _{ab}+{\displaystyle \frac{1}{8}}\sigma ^2(V^2)_{s_as_b}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{f_m^2V_{s_am}V_{ms_b}}{iE+iE_s+i(E_mE_s)f_m}}.`$ (17b)
Corresponding to the changes of variable that have been made, the inversion formulas become
$`E[C_{s_a}(t)]`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{exp}(iE_st){\displaystyle _{iϵ\mathrm{}}^{iϵ+\mathrm{}}}𝑑E\mathrm{exp}(iEt)g_{s_a}(iE+iE_s),`$
$`E[C_{m\{s_a\}}(t)]`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{exp}(iE_mt){\displaystyle _{iϵ\mathrm{}}^{iϵ+\mathrm{}}}𝑑E\mathrm{exp}(iEt)g_m(iE+iE_m).`$ (17c)
Inspecting the equation for the kernel $`K_{ab}`$, we see that apart from order $`V^2`$ terms it is a diagonal matrix $`(iE+iE_s)\delta _{ab}`$. Hence the solution $`g_{s_b}(iE+iE_s)`$, on the inversion contour of integration, will be appreciable only in the vicinity of $`E=iϵ+E_s`$, that is, only near energy shell. This motivates the Weisskopf-Wigner approximation of replacing $`E`$ in the denominator of the final term in $`K_{ab}`$ by $`iϵ+E_s`$, with the result that $`K_{ab}`$ then becomes a linear function of $`E`$. Before making this approximation, the kernel $`K_{ab}`$ has a non-trivial dependence on the stochasticity parameter $`\sigma `$. However, after making the Weisskopf-Wigner approximation, this $`\sigma `$ dependence completely cancels:

$`{\displaystyle \frac{1}{8}}\sigma ^2(V^2)_{s_as_b}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{f_m^2V_{s_am}V_{ms_b}}{iE+iE_s+i(E_mE_s)f_m}}`$
$``$ $`{\displaystyle \frac{1}{8}}\sigma ^2(V^2)_{s_as_b}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{f_m^2V_{s_am}V_{ms_b}}{ϵ+i(E_mE_s)f_m}}`$
$`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(V^2)_{s_as_b}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{f_mV_{s_am}V_{ms_b}}{ϵ+i(E_mE_s)}}`$
$`=`$ $`{\displaystyle \frac{1}{8}}\sigma ^2(V^2)_{s_as_b}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{[1(i/8)\sigma ^2(E_mE_s)]V_{s_am}V_{ms_b}}{ϵ+i(E_mE_s)}}`$
$`=`$ $`{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{iV_{s_am}V_{ms_b}}{E_mE_siϵ}},`$ (18a)
where in the final step we have made use of the condition of Eq. (9a). Thus in the Weisskopf-Wigner approximation, the kernel $`K_{ab}`$ appearing in Eq. (17b) simplifies to
$`K_{ab}`$ $`=`$ $`(iE+iE_s)\delta _{ab}+{\displaystyle \underset{m\{s_a\}}{}}{\displaystyle \frac{iV_{s_am}V_{ms_b}}{E_mE_siϵ}}`$
$`=`$ $`(iE+iE_s)\delta _{ab}+iM_{ab}+{\displaystyle \frac{1}{2}}\mathrm{\Gamma }_{ab},`$
$`M_{ab}`$ $`=`$ $`{\displaystyle \underset{m\{s_a\}}{}}P{\displaystyle \frac{V_{s_am}V_{ms_b}}{E_sE_m}},`$
$`\mathrm{\Gamma }_{ab}`$ $`=`$ $`2\pi {\displaystyle \underset{m\{s_a\}}{}}V_{s_am}V_{ms_b}\delta (E_mE_s),`$ (18b)
with $`P`$ in the definition of the “mass matrix” $`M_{ab}`$ the principal value. These are the same as the formulas for the kernel in the absence of the stochastic terms in the Schrödinger equation. Thus, in the Weisskopf-Wigner approximation, the solution for $`E[C_{s_a}(t)]`$ is unmodified by the stochastic effects, and hence the Lorentzian line profile and the decay rate of the state are unaffected by the $`\sigma `$ terms.

The solution for $`E[C_m(t)]`$ with $`m\{s_a\}`$ does retain a dependence on the stochastic parameter. To study this, let us specialize to the case $`D=1`$ of a non-degenerate initial state. The expression in Eq. (18b) for the kernel now becomes the $`1\times 1`$ matrix

$$K(E)=iE+iE_s+iM+\frac{1}{2}\mathrm{\Gamma },$$
(19a)
with $`M`$ and $`\mathrm{\Gamma }`$ real numbers given by
$`M`$ $`=`$ $`{\displaystyle \underset{ms}{}}P{\displaystyle \frac{V_{sm}V_{ms}}{E_sE_m}},`$
$`\mathrm{\Gamma }`$ $`=`$ $`2\pi {\displaystyle \underset{ms}{}}V_{sm}V_{ms}\delta (E_mE_s).`$ (19b)
Thus, Eq. (17b) has the immediate solution
$$g_s(iE+iE_s)=K(E)^1,$$
(19c)
which when substituted into Eq. (17a) yields
$$g_m(iE+iE_m)=[iE+iE_m+\frac{1}{8}\sigma ^2(E_mE_s)^2]^1V_{ms}f_m(EE_sM+\frac{i}{2}\mathrm{\Gamma })^1.$$
(19d)
Substituting these equations into the inversion formulas of Eqs. (17c), and doing elementary contour integrations, we find

$`E[C_s(t)]`$ $`=`$ $`\mathrm{exp}(iMt{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t),`$ (20a)
$`E[C_{ms}(t)]`$ $`=`$ $`{\displaystyle \frac{V_{ms}}{E_sE_m+M\frac{i}{2}\mathrm{\Gamma }}}\left(\mathrm{exp}[i(E_mE_sM)t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{exp}[{\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2t]\right).`$
¿From Eq. (11b) we thus get
$$E[|C_s(t)|^2]=\mathrm{exp}(\mathrm{\Gamma }t),$$
(20b)
which identifies $`\mathrm{\Gamma }`$ as the transition rate per unit time out of the initial state. Finally, substituting Eq. (20a) into Eq. (12b), simplifying to leading order in $`V`$, and integrating with respect to $`t`$, we get
$`E[|C_{ms}(t)|^2]`$ $`=`$ $`{\displaystyle \frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2}}(\mathrm{exp}(\mathrm{\Gamma }t)+1`$ (20c)
$`2\mathrm{exp}[{\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{cos}[(E_sE_m+M)t]).`$
This completes our solution for the expectations of the coefficients, and their squared magnitudes, in the case of a non-degenerate initial state. We see that after a time $`t`$ large compared with the lifetime $`\mathrm{\Gamma }^1`$, we obtain
$$E[|C_{ms}(\mathrm{})|^2]=\frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2},$$
(20d)
exhibiting the standard Lorentzian profile with no dependence on the stochasticity parameter $`\sigma `$.

## VII Small Time and Golden Rule Approximations

Let us now study the behavior of Eq. (20c) for small and large values of the time $`t`$. Since within our approximations we have $`\sigma ^2(E_mE_s)^2\sigma ^2[(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2]`$, we can rewrite Eq. (20c) as

$`E[|C_{ms}(t)|^2]`$ $`=`$ $`{\displaystyle \frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2}}(\mathrm{exp}(\mathrm{\Gamma }t)+1`$
$`2\mathrm{exp}[{\displaystyle \frac{1}{8}}\sigma ^2((E_sE_m+M)^2+{\displaystyle \frac{1}{4}}\mathrm{\Gamma }^2)t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{cos}[(E_sE_m+M)t]).`$
In the limit as $`t0`$, we can develop the exponential and cosine functions in power series expansions, giving
$$E[|C_{ms}(t)|^2]|V_{ms}|^2\left(\frac{1}{4}\sigma ^2t+t^2+O(t^2\sigma ^4(E_sE_m)^2)+O(t^3)\right).$$
(21b)
Thus the leading small time behavior of the summed expected probability in the decay channels is
$$\underset{ms}{}|V_{ms}|^2\frac{1}{4}\sigma ^2t=(V^2)_{ss}\frac{1}{4}\sigma ^2t,$$
(21c)
where in evaluating the sum we have employed the condition of Eq. (9a). We shall verify this result by another method in Sec. IX, where we discuss its implications for the quantum Zeno effect, and in Sec. X shall apply it to estimating bounds on $`\sigma `$.

Let us next consider the large time behavior implied by Eq. (21). Once $`t`$ is large enough so that $`|(E_mE_s)t|`$ is large for all energies $`E_m`$ not infinitesimally close to $`E_s`$, we can evaluate the summed expected probability in the decay channels by making the “golden rule” approximation multi . This approximation treats the factors multiplying $`|V_{ms}|^2`$ in Eq. (21), which are sharply peaked around $`E_m=E_s`$, as if they were equal to a Dirac delta function of strength given by the integral of these factors over energy. We then have

$`{\displaystyle \underset{ms}{}}E[|C_{ms}(t)|^2]`$ $`=`$ $`{\displaystyle \underset{ms}{}}{\displaystyle \frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2}}(\mathrm{exp}(\mathrm{\Gamma }t)+1`$ (22a)
$``$ $`2\mathrm{exp}[{\displaystyle \frac{1}{8}}\sigma ^2((E_sE_m+M)^2+{\displaystyle \frac{1}{4}}\mathrm{\Gamma }^2)t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{cos}[(E_sE_m+M)t])`$
$``$ $`{\displaystyle \underset{ms}{}}|V_{ms}|^2\delta (E_mE_s)`$
$`\times `$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}d(\mathrm{\Delta }E){\displaystyle \frac{\mathrm{exp}(\mathrm{\Gamma }t)+12\mathrm{exp}[\frac{1}{8}\sigma ^2\left((\mathrm{\Delta }E)^2+\frac{1}{4}\mathrm{\Gamma }^2\right)t\frac{1}{2}\mathrm{\Gamma }t]\mathrm{cos}[(\mathrm{\Delta }E)t]}{(\mathrm{\Delta }E)^2+\frac{1}{4}\mathrm{\Gamma }^2}}`$
$`=`$ $`{\displaystyle \underset{ms}{}}|V_{ms}|^2\delta (E_mE_s)tF[\sigma ^2/(8t),t]`$
$`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }t}{2\pi }}F[\sigma ^2/(8t),t],`$
with $`\mathrm{\Gamma }`$ as given in Eq. (19b) and with the function $`F[A,t]`$ defined by
$$F[A,t]=_{\mathrm{}}^{\mathrm{}}𝑑u\frac{\mathrm{exp}(\mathrm{\Gamma }t)+12\mathrm{exp}[A(u^2+\frac{1}{4}\mathrm{\Gamma }^2t^2)\frac{1}{2}\mathrm{\Gamma }t]\mathrm{cos}u}{u^2+\frac{1}{4}\mathrm{\Gamma }^2t^2}.$$
(22b)
To evaluate $`F[A,t]`$ we note that gradsh

$$F[0,t]=_{\mathrm{}}^{\mathrm{}}𝑑u\frac{\mathrm{exp}(\mathrm{\Gamma }t)+12\mathrm{exp}(\frac{1}{2}\mathrm{\Gamma }t)\mathrm{cos}u}{u^2+\frac{1}{4}\mathrm{\Gamma }^2t^2}=\frac{2\pi }{\mathrm{\Gamma }t}[1\mathrm{exp}(\mathrm{\Gamma }t)],$$
(23a)
and gradsh
$`{\displaystyle \frac{F[A,t]}{A}}`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑u2\mathrm{exp}[A(u^2+{\displaystyle \frac{1}{4}}\mathrm{\Gamma }^2t^2){\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{cos}u`$ (23b)
$`=`$ $`2\pi ^{1/2}\mathrm{exp}({\displaystyle \frac{1}{2}}\mathrm{\Gamma }t{\displaystyle \frac{1}{4}}A\mathrm{\Gamma }^2t^2)A^{1/2}\mathrm{exp}[1/(4A)].`$
Thus, integrating Eq. (23b) with respect to $`A`$ we get
$$F[A,t]=\frac{2\pi }{\mathrm{\Gamma }t}[1\mathrm{exp}(\mathrm{\Gamma }t)]+C[A,t],$$
(23c)
with the correction term $`C[A,t]`$ given by

$$C[A,t]=4\pi ^{1/2}\mathrm{exp}(\frac{1}{2}\mathrm{\Gamma }t)_0^{A^{1/2}}𝑑v\mathrm{exp}[\frac{1}{4}(v^2\mathrm{\Gamma }^2t^2+1/v^2)].$$
(24a)
Since the exponentials of negative arguments in Eq. (24a) are bounded by their maxima over the range of integration, we have
$$|C[A,t]|<4\pi ^{1/2}A^{1/2}\mathrm{exp}[1/(4A)]=2\sigma (\frac{\pi }{2t})^{1/2}\mathrm{exp}(2t/\sigma ^2).$$
(24b)
So when $`\mathrm{\Gamma }t`$ is of order unity, the contribution of the correction term $`C[A,t]`$ is of order $`c\sigma \mathrm{\Gamma }^{1/2}\mathrm{exp}[2/(\sigma ^2\mathrm{\Gamma })]`$ $`c^{}\sigma V\mathrm{exp}[c^{\prime \prime }/(\sigma V)^2]`$, with $`c,c^{},c^{\prime \prime }`$ constants, which is exponentially small and can be neglected in our approximation scheme. Thus we are justified in approximating
$$F[A,t]F[0,t]=\frac{2\pi }{\mathrm{\Gamma }t}[1\mathrm{exp}(\mathrm{\Gamma }t)],$$
(24c)
which when substituted back into Eq. (22a) gives
$$\underset{ms}{}E[|C_{ms}(t)|^2]=1\mathrm{exp}(\mathrm{\Gamma }t)=1|C_s(t)|^2,$$
(24d)
verifying that the approximations used in our calculation are consistent with maintenance of the unitarity sum rule (the unit state vector normalization condition).

## VIII Solution to the Stochastic Equation for $`C_{ms}`$

Since we see from Eqs. (20a) and (20c) that $`E[|C_{ms}(t)|^2`$ differs from $`|E[C_{ms}(t)]|^2`$, the stochastic fluctuations in $`C_{ms}(t)`$ are evidently playing a role. Let us now demonstrate this directly by solving the stochastic differential equation for $`C_{ms}(t)`$. Specializing to the case of a non-degenerate initial state, approximating $`C_s(t)E[C_s(t)]`$, and using Eq. (20a) for $`E[C_s(t)]`$, Eq. (10d) becomes

$`dC_m(t)`$ $`=`$ $`(\alpha _m^{(1)}dW_t+\alpha _m^{(2)}dt)C_m(t)`$
$`+`$ $`\mathrm{exp}[i(E_mE_sM)t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t](\gamma _{ms}^{(1)}dW_t+\gamma _{ms}^{(2)}dt),`$
$`\alpha _m^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma (E_mE_s),\alpha _m^{(2)}={\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2={\displaystyle \frac{1}{2}}(\alpha _m^{(1)})^2,`$
$`\gamma _{ms}^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma V_{ms},\gamma _{ms}^{(2)}iV_{ms}f_m.`$ (25a)
For general values of the coefficients $`\alpha _m^{(1,2)}`$ and $`\gamma _{ms}^{(1,2)}`$, Eq. (25a) can be integrated by using Eqs. (2a-c) to find a stochastic integrating factor for the $`C_m`$ terms (see the Appendix), with the result
$`C_m(t)`$ $`=`$ $`\mathrm{exp}[\alpha _m^{(1)}W_t(\alpha _m^{(1)})^2t]{\displaystyle _0^t}\mathrm{exp}[i(E_mE_sM)u{\displaystyle \frac{1}{2}}\mathrm{\Gamma }u\alpha _m^{(1)}W_u+(\alpha _m^{(1)})^2u]`$ (25b)
$`\times `$ $`[\gamma _{ms}^{(1)}dW_u+(\gamma _{ms}^{(2)}\alpha _m^{(1)}\gamma _{ms}^{(1)})du].`$
Using Eqs. (2a-c), it is easy to verify directly that Eq. (25b) solves Eq. (25a). If we now examine Eq. (25b) more closely, using the specific expressions for the coefficients $`\alpha _m^{(1,2)}`$ and $`\gamma _{ms}^{(1,2)}`$ given in Eq. (25a), we find that within the approximation of neglecting terms of relative order $`\sigma V^2`$, the integrand in Eq. (25b) is an exact stochastic differential. Thus the integration can be carried out explicitly (see the Appendix), with the result
$`C_{ms}(t)`$ $`=`$ $`{\displaystyle \frac{V_{ms}}{E_sE_m+M\frac{i}{2}\mathrm{\Gamma }}}`$
$`\times `$ $`\left(\mathrm{exp}[i(E_mE_sM)t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{exp}[{\displaystyle \frac{1}{2}}\sigma (E_mE_s)W_t{\displaystyle \frac{1}{4}}\sigma ^2(E_mE_s)^2t]\right).`$
This expression can be easily verified, by use of Eqs. (2a-c), to be the solution to Eq. (25a) (up to a residual error of relative order $`\sigma V^2`$). Using Eq. (3c) to take the expectation of Eq. (25), we recover the result of Eq. (20a). From Eq. (25) we find an explicit formula for $`|C_{ms}(t)|^2`$,
$`|C_{ms}(t)|^2`$ $`=`$ $`{\displaystyle \frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2}}`$
$`\times `$ $`(\mathrm{exp}(\mathrm{\Gamma }t)+\mathrm{exp}[\sigma (E_mE_s)W_t{\displaystyle \frac{1}{2}}\sigma ^2(E_mE_s)^2t]`$
$``$ $`2\mathrm{exp}[{\displaystyle \frac{1}{2}}\sigma (E_mE_s)W_t{\displaystyle \frac{1}{4}}\sigma ^2(E_mE_s)^2t{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{cos}(E_mE_sM)t).`$
Again using Eq. (3c) to take the expectation of this formula, we recover the result of Eq. (20c).

## IX Stochastic Suppression of the Quantum Zeno Effect

In Eqs. (21b,c) we saw that the leading small time behavior of the summed expected probability in the decay channels is

$$(V^2)_{ss}\frac{1}{4}\sigma ^2t,$$
(26)
rather than the result $`(V^2)_{ss}t^2`$ that would hold for vanishing $`\sigma `$. As a result, $`E[|C_s(t)|^2]1`$ vanishes linearly in $`t`$ for nonzero $`\sigma `$, rather than quadratically in $`t`$ as for the unmodified Schrödinger equation. Since the quadratic vanishing of $`|C_s(t)|^21`$ in standard quantum mechanics is the origin of the quantum Zeno effect misra , we conclude that in the energy driven stochastic Schrödinger equation, the quantum Zeno effect is suppressed.

Let us verify this directly from the stochastic differential equation of Eq. (1a), in analogy with the direct calculation anan  of $`|C_s(t)|^21`$ for small times for the ordinary Schrödinger equation. Applying the Itô rule of Eq. (2a), we have

$$d|s(0)|s(t)|^2|_{t=0}=s(0)|d|s(t)|_{t=0}+s(0)|d|s(t)^{}|_{t=0}+s(0)|d|s(t)|_{t=0}s(0)|d|s(t)^{}|_{t=0}.$$
(27a)
From Eq. (1a) we have
$$d|s(t)=iH|s(t)dt\frac{1}{8}\sigma ^2(Hs(t)|H|s(t))^2|s(t)dt+\frac{1}{2}\sigma (Hs(t)|H|s(t))|s(t)dW_t,$$
(27b)
and so setting $`t=0`$ and projecting on $`s(0)|`$ gives
$$s(0)|d|s(t)|_{t=0}=iH_sdt\frac{1}{8}\sigma ^2(HH_s)^2_sdt,$$
(27c)
with $`H^n_s=s(0)|H^n|s(0)`$. Substituting Eq. (27c) into Eq. (27a), we thus get the first term in the small $`t`$ expansion of $`|s(0)|s(t)|^21`$,

$$|s(0)|s(t)|^2=1\frac{1}{4}\sigma ^2(HH_s)^2_st+O(t^2).$$
(28a)
This equation gives a general formula for the stochastic suppression of the quantum Zeno effect, independent of any assumptions about the potential. When the general form anan  of the order $`t^2`$ term coming from the standard Schrödinger evolution is included, Eq. (28a) becomes
$$|s(0)|s(t)|^2=1(HH_s)^2_s\left(\frac{1}{4}\sigma ^2t+t^2\right)+O(\sigma ^4t^2)+O(t^3);$$
(28b)
in other words, the first two terms in the small $`t`$ expansion are governed to leading order in $`\sigma `$ by the initial state energy variance. When the potential is assumed to obey Eq. (9a), we have
$`H_s`$ $`=`$ $`E_s+V_{ss}=E_s,`$
$`H^2_s`$ $`=`$ $`E_s^2+2E_sV_{ss}+(V^2)_{ss},`$
$`(HH_s)^2_s`$ $`=`$ $`H^2_sH_s^2=(V^2)_{ss},`$ (28c)
and so Eq. (28b) becomes
$$|s(0)|s(t)|^2=1(V^2)_{ss}\left(\frac{1}{4}\sigma ^2t+t^2\right)+O(\sigma ^4t^2)+O(t^3),$$
(28d)
in agreement with the result of Eqs. (21a,b) and the unitarity sum rule.

## X Discussion and Estimates of Bounds on $`\sigma `$

We have seen that to leading order in the perturbing potential, the stochastic terms governed by $`\sigma `$ do not affect either the Lorentzian line profile or the transition rate per unit time as evaluated in the Weisskopf-Wigner approximation, but only produce a change in the short time transient behavior of the transition probabilities from the initial state. This is a direct result of the fact that the energy-driven stochastic Schrödinger equation is energy conserving. On dimensional grounds, the transition rate per unit time $`\mathrm{\Gamma }`$ could contain, in addition to the usual terms of the form $`\delta (E_sE_m)|V_{ms}|^2`$, a term of the form $`\sigma ^2(V^2)_{ss}`$. However, this additional term is not energy conserving, and as a result we have seen that its coefficient precisely cancels to zero in the Weisskopf-Wigner approximation.

Because the transition rate per unit time and Lorentzian line shape are unaffected by $`\sigma `$, bounds on $`\sigma `$ from particle decays result only from experiments in which a metastable system is monitored as function of time from a known time (or vertex location) of formation. According to Eqs. (28a-d), for small times the effective transition rate per unit time is

$$\mathrm{\Gamma }_R=\frac{1}{4}\sigma ^2(\mathrm{\Delta }E)^2=\frac{1}{4}\sigma ^2(V^2)_{ss},$$
(29a)
with $`(\mathrm{\Delta }E)^2=(HH_s)^2_s`$ the initial state energy variance. This can be interpreted as an early time decay rate coming from spontaneous reduction induced by the stochastic fluctuation terms, in agreement with the estimate $`\mathrm{\Gamma }_R\sigma ^2(\mathrm{\Delta }E)^2`$ used in earlier discussions gisin ; adler . In order for the rate of Eq. (29a) to not lead to pronounced early time deviations from the observed decay rate $`\mathrm{\Gamma }`$, we must have
$$\mathrm{\Gamma }_R<\mathrm{\Gamma },$$
(29b)
which writing $`\sigma ^2=M_\sigma ^1`$ implies the bound

$$M_\sigma >\frac{(V^2)_{ss}}{4\mathrm{\Gamma }}=\frac{\underset{ms}{}|V_{sm}|^2}{8\pi _{ms}|V_{sm}|^2\delta (E_mE_s)}\frac{E_D}{8\pi },$$
(30)
with $`E_D`$ defining an energy characteristic of the decay process. In a particle physics context, a first guess would be to estimate $`E_D`$ as being of order the mass of the decaying particle. The most massive decays for which $`\mathrm{\Gamma }`$ has been measured by tracking a metastable system from the point of formation appear to be $`\pi ^0\gamma \gamma `$ decay, with an initial mass order 140 MeV, and charmed meson decays, with an initial mass of around 2 GeV. Estimating $`E_D`$ in Eq. (30) as the decaying particle mass, these give respective bounds on $`M_\sigma `$ of order 6 MeV and 80 MeV, respectively. If $`M_\sigma `$ were significantly larger than these bounds, one would have observed anomalous accumulations of decay events close to the production vertex, as a result of decays induced by spontaneous reduction. For comparison, the observation of coherent superpositions of energy eigenstates in the neutrino, $`K`$-meson and $`B`$-meson systems gives bounds sladler , respectively, of $`M_\sigma >10^{20}`$GeV, $`M_\sigma >2\times 10^{15}`$GeV, and $`M_\sigma >2\times 10^{13}`$GeV.

Thus the charmed meson decay bound on $`M_\sigma `$ represents a significant improvement over the coherent oscillation bounds. However, it is still smaller than the Planck mass, which is very likely the expected value of $`M_\sigma `$, by a factor of $`10^{20}`$! We conclude that the theory of decaying states in the energy-driven stochastic Schrödinger equation places only very weak empirical bounds on the magnitude of the stochasticity parameter $`\sigma `$.

We leave for future study two issues that can be addressed within the general framework established here. The first is an analysis of the nature of the transition between the short-time regime with decay rate $`\mathrm{\Gamma }_R`$, and the exponential decay regime with decay rate $`\mathrm{\Gamma }`$. This is governed by the solution of Eqs. (17a,b) before making the Weisskopf-Wigner approximation of replacing $`E`$ in the denominator of the final term in $`K_{ab}`$ by $`E_s`$. The second is an analysis of the magnitude of the energy $`E_D`$ defined by Eq. (30), for various dynamical models of the decay process, as reflected in the energy spectrum of the unperturbed states $`|m`$ and in the magnitudes of the decay-inducing matrix elements $`V_{sm}`$.

## XI Acknowledgments

This work was supported in part by the Department of Energy under Grant #DE–FG02–90ER40542. I wish to acknowledge the hospitality of the Aspen Center for Physics, where most of the calculations reported here were done. I also wish to thank Angelo Bassi, Dorje Brody, Todd Brun, and especially Lane Hughston, for informative discussions about the properties of the Itô calculus, Todd Brun for a helpful discussion about extracting bounds from decay observations, Larry Horwitz for a discussion several years ago about virtues of the Laplace transform, and Edward Witten for raising the issue of empirical bounds on the stochasticity parameter.

## XII Added Note

Lajos Diósi ldiosi  has pointed out an elegant stochastic-theoretic technique that allows the main physical results of this paper to be derived in a few lines, starting from the standard quantum mechanical results that hold when the stochasticity parameter $`\sigma `$ is zero. Diósi makes three principal observations. The first is that the quantities of direct physical interest, as pointed out in Sec. V, are the expectations $`E[|C_m(t)|^2]`$ of the squared magnitudes of the perturbation coefficients. Since according to Eq. (6b) we have

$$|C_m(t)|^2=|m|\psi (t)|^2=m|\psi (t)\psi (t)|m=m|\rho (t)|m,$$
(31a)
with $`\rho (t)`$ the density matrix
$$\rho (t)=|\psi (t)\psi (t)|,$$
(31b)
to calculate $`E[|C_m(t)|^2`$ it suffices to know $`E[\rho (t)]`$, in other words
$$E[|C_m(t)|^2]=m|E[\rho (t)]|m=\mathrm{Tr}(|mm|)E[\rho (t)].$$
(31c)
Since the dynamics of $`E[\rho (t)]`$ is governed by the Lindblad-type equation of Eq. (5), to calculate the physically relevant expectations it thus suffices to solve the dynamical problem specified by Eq. (5), supplemented by the initial condition
$$E[\rho (0)]=\rho (0)=|s_As_A|.$$
(31d)
Diósi’s second observation is that the dynamical problem specified by Eq. (5), with the initial condition of Eq. (31d), can be compactly solved by a simple stochastic trick. The trick uses the fact noise  that there is a second stochastic Schrödinger equation, simpler in structure than that of Eq. (1a), which also leads to Eq. (5) as the equation for the evolution of the stochastic expectation of its density matrix. To see this, consider the stochastic Schrödinger equation

$$d|\psi =iH|\psi dt\frac{1}{8}\sigma ^2H^2|\psi dt+\frac{1}{2}i\sigma H|\psi dW_t.$$
(32a)
This equation differs from that of Eq. (1a) in having an imaginary noise term, with operator coefficient $`H`$, instead of a real noise term with operator coefficient $`HH`$. A simple calculation, using the Itô calculus rules of Eq. (1b), shows that Eq. (32a) also leads to preservation of the norm of the state $`|\psi `$, and leads to the density matrix evolution equation
$$d\rho =i[\rho ,H]dt\frac{1}{8}\sigma ^2[H,[H,\rho ]]dt+\frac{1}{2}i\sigma [H,\rho ]dW_t,$$
(32b)
which has the stochastic expectation
$$dE[\rho ]=i[E[\rho ],H]dt\frac{1}{8}\sigma ^2[H,[H,E[\rho ]]]dt,$$
(32c)
which is identical to Eq. (5). Hence the imaginary noise equation of Eq. (32a) will lead to the same results for the physical quantities $`E[|C_m(t)|^2]`$ as the real noise equation of Eq. (1a), even though the stochastic details of the two processes differ!

Diósi’s third observation is the fact that Eq. (32a) can be immediately formally integrated to give

$$|\psi (t)=\mathrm{exp}[iH(t\frac{1}{2}\sigma W_t)]|\psi (0,$$
(33a)
as can be readily ascertained by use of Eq. (2c) with the choice
$$\alpha =\frac{1}{2}i\sigma H.$$
(33b)
Combining this observation with the first two, then leads to a very simple rule for calculating the stochastic modifications of decay processes governed by Eq. (1a). Let $`E[|C_m^\sigma (t)|^2]`$ be the quantities of physical interest, viewed as functions of $`\sigma `$ as well as of $`t`$, so that $`E[|C_m^0(t)|^2]=|C_m^0(t)|^2`$ are their values as calculated from the standard Schrödinger evolution with no stochasticity. Then Eqs. (31a) through (33b) imply the simple relation

$$E[|C_m^\sigma (t)|^2]=E[|C_m^0(t\frac{1}{2}\sigma W_t)|^2],$$
(34a)
between the probabilities calculated in the standard Schrödinger analysis, and the stochastic expectations of the probabilities as calculated from Eq. (1a). The recipe is simply this: take the known expressions for the probabilities calculated in standard quantum mechanics, replace $`t`$ by $`t\frac{1}{2}\sigma W_t`$, and take the stochastic expectation. The needed stochastic expectations of powers of $`W_t`$ can all be read off from the expansion of Eq. (3c) in powers of $`\alpha `$,
$$E[W_t]=0,E[W_t^2]=t,E[W_t^3]=0,E[W_t^4]=3t^2,\mathrm{}.$$
(34b)
Let us now apply Diósi’s observations to rederive the principal results found above for the stochastic analog of the Weisskopf-Wigner analysis. First, let us consider the short time behavior of the survival probability given in Eq. (28b). The standard answer when $`\sigma =0`$, which gives the quantum Zeno effect, is

$$|s(0)|s(t)|^2=1(HH_s)^2_st^2+O(t^3).$$
(35a)
Following the recipe, we have
$$E[(t\frac{1}{2}\sigma W_t)^2]=E[t^2t\sigma W_t+\frac{1}{4}\sigma ^2W_t^2]=t^2+\frac{1}{4}\sigma ^2t.$$
(35b)
On substitution into Eq. (35a) this gives for general $`\sigma `$
$$|s(0)|s(t)|^2=1(HH_s)^2_s(t^2+\frac{1}{4}\sigma ^2t)+\mathrm{},$$
(35c)
in agreement with the result for the stochastic modification of the quantum Zeno effect given in Eq. (28b).

Next let us apply the recipe to the formula for the initial state survival probability obtained using the Weisskopf-Wigner approximation, which is valid for times $`t`$ that are not too small (and also not too large). The standard analysis gives

$$|C_s^0(t)|^2=\mathrm{exp}(\mathrm{\Gamma }t),$$
(36a)
with $`\mathrm{\Gamma }`$ the Golden Rule decay rate of Eq. (19b). Replacing $`t`$ by $`t\frac{1}{2}\sigma W_t`$ and using Eq. (3c) to take the stochastic expectation, we get as the exact formula for the stochastic modification of the Weisskopf-Wigner approximation
$$E[|C_s^\sigma (t)|^2]=\mathrm{exp}[\mathrm{\Gamma }(1\frac{1}{8}\sigma ^2\mathrm{\Gamma })t],$$
(36b)
which reduces, when the correction term of relative order $`\sigma ^2\mathrm{\Gamma }`$ is neglected, to the answer found in Eq. (20b). Since $`\sigma ^2\mathrm{\Gamma }O(\sigma ^2V^2)`$, we see that the calculation of Secs. 4-6 above did not succeed in keeping all terms of order $`\sigma ^2V^2`$, and in fact there is a small stochastic correction to the decay rate, with the corrected decay rate given by
$$\mathrm{\Gamma }^\sigma =\mathrm{\Gamma }(1\frac{1}{8}\sigma ^2\mathrm{\Gamma }).$$
(36c)
However, writing $`\sigma ^2=M_\sigma ^1`$ as in Sec. X, as long as $`M_\sigma >E_s`$ this correction is not significant within the Weisskopf-Wigner approximation scheme, which treats the line width $`\mathrm{\Gamma }`$ as a small quantity relative to $`E_s`$.

Finally, let us apply the recipe to the formula giving the probability for a transition to the state $`|m`$. The standard Weisskopf-Wigner approximation result for this is given by Eq. (20c) with $`\sigma =0`$,

$`|C_{ms}^0(t)|^2`$ $`=`$ $`{\displaystyle \frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2}}(\mathrm{exp}(\mathrm{\Gamma }t)+1`$ (37a)
$`2\mathrm{exp}[{\displaystyle \frac{1}{2}}\mathrm{\Gamma }t]\mathrm{cos}[(E_sE_m+M)t]).`$
Applying the recipe, and again using Eq. (3c) to evaluate the needed expectations, we get the exact stochastic extension of Eq. (37a),
$`E[|C_{ms}^\sigma (t)|^2]`$ $`=`$ $`{\displaystyle \frac{|V_{ms}|^2}{(E_sE_m+M)^2+\frac{1}{4}\mathrm{\Gamma }^2}}(\mathrm{exp}[\mathrm{\Gamma }(1{\displaystyle \frac{1}{8}}\sigma ^2\mathrm{\Gamma })t]+1`$ (37b)
$`2\mathrm{exp}[{\displaystyle \frac{1}{2}}\mathrm{\Gamma }(1{\displaystyle \frac{1}{16}}\sigma ^2\mathrm{\Gamma })t{\displaystyle \frac{1}{8}}\sigma ^2(E_sE_m+M)^2t]`$
$`\times `$ $`\mathrm{cos}[(E_sE_m+M)(1{\displaystyle \frac{1}{8}}\sigma ^2\mathrm{\Gamma })t]).`$
Again, when simplified to leading order in $`V`$, this gives the result of Eq. (20c) above. However, even before dropping nonleading terms in $`V`$, we see that Eq. (37b) implies the Lorentzian formula of Eq. (20d) in the large time limit.

From the above exposition, we see that Diósi’s observations not only greatly simplify the calculation of the physically relevant quantities, but also give results that are completely independent of the assumption of Eq. (9a) that was used to linearize the stochastic equation. (This is something that one might have already suspected from the fact that Eq. (28b) is more general than Eq. (28d).) Thus, the only approximations that are needed to get stochastic results are those that are used in the standard, non-stochastic quantum mechanical analysis. Moreover, the “miraculous” cancellation of the $`\sigma ^2`$ terms in the Weisskopf-Wigner approximation to the mass and decay matrices, exhibited above in Eq. (18a), is given a deeper explanation. There is an extensive literature facchi  discussing the decay problem without making the Weisskopf-Wigner approximation (i.e., without replacing $`E`$ by $`E_s`$ in the order $`V^2`$ terms of the Laplace transform kernel $`K_{ab}`$), and these discussions can all be converted to results for $`E[|C_m^\sigma (t)|^2]`$ in the stochastic case, by using the recipe of replacing $`t`$ by $`t\frac{1}{2}\sigma W_t`$ in the corresponding formula for $`|C_m^0(t)|^2`$ and taking a stochastic average over $`W_t`$.

The relation of Eq. (34a) between stochastic and standard quantum mechanical probabilities can be applied to other problems as well. For example, the density matrix of a two-level system can be represented in the form

$$\rho =\frac{1}{2}(1\stackrel{}{R}\stackrel{}{\tau }),$$
(38a)
with $`\stackrel{}{\tau }=(\tau _1,\tau _2,\tau _3)`$ the standard Pauli matrices, and with $`\stackrel{}{R}=(R_1,R_2,R_3)`$ a vector summarizing the structure of the traceless part of the density matrix. The standard, $`\sigma =0`$ Schrödinger equation describing Rabi oscillations of the two-level system under the influence of an applied field oscillating at the frequency of the level separation (in co-rotating coordinates, neglecting the counter-rotating field component) gives for the equation of motion iirabi  of the vector $`\stackrel{}{R}|_{\sigma =0}\stackrel{}{R}^0`$,
$$\frac{d\stackrel{}{R}^0}{dt}=\stackrel{}{\omega }\times \stackrel{}{R}^0,$$
(38b)
with $`|\stackrel{}{\omega }|=\mathrm{\Omega }`$ the angular frequency of precession of $`\stackrel{}{R}^0`$. The probabilities for finding the system in the upper and lower levels are given, as a function of time, by
$$P_\pm ^0(t)=\frac{1}{2}[1\pm R_3^0(t)].$$
(38c)
Since the general solution of Eq. (38b) has the form

$$\stackrel{}{R}^0(t)=\stackrel{}{V}_1\mathrm{cos}\mathrm{\Omega }t+\stackrel{}{V}_2\mathrm{sin}\mathrm{\Omega }t,$$
(39a)
with $`\stackrel{}{V}_{1,2}`$ fixed vectors that depend on the initial state and the structure of the Hamiltonian, and since
$`E[\mathrm{cos}\mathrm{\Omega }(t{\displaystyle \frac{1}{2}}\sigma W_t)]`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{1}{8}}\mathrm{\Omega }^2\sigma ^2t)\mathrm{cos}\mathrm{\Omega }t,`$
$`E[\mathrm{sin}\mathrm{\Omega }(t{\displaystyle \frac{1}{2}}\sigma W_t)]`$ $`=`$ $`\mathrm{exp}({\displaystyle \frac{1}{8}}\mathrm{\Omega }^2\sigma ^2t)\mathrm{sin}\mathrm{\Omega }t,`$ (39b)
we have under the stochastic evolution of Eq. (1a)
$$E[\stackrel{}{R}^\sigma (t)]=\mathrm{exp}(\frac{1}{8}\mathrm{\Omega }^2\sigma ^2t)\stackrel{}{R}^0(t).$$
(39c)
By Eq. (38c), this gives for the expected probabilities when the system evolves under the stochastic Schrödinger equation,
$$\frac{1}{2}E[P_\pm ^\sigma (t)]=\mathrm{exp}(\frac{1}{8}\mathrm{\Omega }^2\sigma ^2t)[\frac{1}{2}P_\pm ^0(t)].$$
(39d)
This can be applied, for example, to the quantum Zeno effect experiment of Itano et. al. itano , who carry out a proposal of Cook cook  to make repeated measurements of a two-level system while the vector $`\stackrel{}{R}`$ is precessing for a time interval $`t=\pi /\mathrm{\Omega }`$, for which the exponential damping factor in Eq. (39c) becomes $`\mathrm{exp}(\frac{1}{8}\pi \mathrm{\Omega }\sigma ^2)`$. Corresponding to the experimental value $`\mathrm{\Omega }=320.7`$ MHz and the fact that probabilities were observed to an accuracy of about $`.02`$ in this experiment, and were found to agree with the standard Schrödinger theory, we get a bound on $`M_\sigma =1/\sigma ^2`$ of $`M_\sigma >2\times 10^{15}`$ GeV, comparable to that obtained from oscillations in the $`K`$-meson system.

*

## Appendix A

We give here the details of the integration of the stochastic differential equation that appears in Sec. VIII. Consider the linear stochastic differential equation

$$dC_t=(A_tdW_t+B_tdt)C_t+P_tdW_t+Q_tdt,$$
(40)
which is to be solved for the unknown stochastic function $`C_t`$ given the known functions $`A_t`$, $`B_t`$, $`P_t`$, and $`Q_t`$. Although we shall proceed as if these known functions were deterministic, in fact all our manipulations and the final solution are unchanged hughston  if the input functions are themselves stochastic. To solve Eq. (40), we transpose the $`C_t`$ term on the right to the left and multiply by a factor $`F_t`$, which is to be determined, giving

$$F_t[dC_t(A_tdW_t+B_tdt)C_t]=F_t[P_tdW_t+Q_tdt].$$
(41)
We now look for an $`F_t`$ which makes the left hand side of Eq. (41) a total differential, up to terms independent of $`C_t`$ that are of the same form as the terms on the right hand side. Making the Ansatz

$$F_t=\mathrm{exp}\left[_0^t(\alpha _udW_u+\beta _udu)\right],$$
(42)
we find by use of Eqs. (2a) and (2c) of the text that

$$d(F_tC_t)=F_t[dC_t+\alpha _tdW_tC_t+(\beta _t+\frac{1}{2}\alpha _t^2)dtC_t+\alpha _tdW_tdC_t],$$
(43)
which on substituting Eq. (40) for the final $`dC_t`$ on the right, and using Eq. (1b) of the text, gives

$$d(F_tC_t)=F_t[dC_t+\alpha _tdW_tC_t+(\beta _t+\frac{1}{2}\alpha _t^2+\alpha _tA_t)dtC_t+\alpha _tP_tdt].$$
(44a)
Hence if we choose
$$\alpha _t=A_t,\beta _t=B_t+\frac{1}{2}A_t^2,$$
(44b)
then Eq. (44a) takes the form

$$d(F_tC_t)=F_t[dC_t(A_tdW_t+B_tdt)C_tA_tP_tdt],$$
(45)
which by use of Eq. (41) becomes

$$d(F_tC_t)=F_t[P_tdW_t+(Q_tA_tP_t)dt].$$
(46)
The dependence on the unknown function $`C_t`$ is now entirely in the form of an exact differential, and so Eq. (46) can be immediately integrated to give

$`C_t`$ $`=`$ $`F_t^1\left[C_0+{\displaystyle _0^t}𝑑uF_u\left(P_udW_u+(Q_uA_uP_u)du\right)\right]`$
$`=`$ $`\mathrm{exp}\left[{\displaystyle _0^t}\left(A_udW_u+(B_u{\displaystyle \frac{1}{2}}A_u^2)du\right)\right]`$
$`\times `$ $`\left(C_0+{\displaystyle _0^t}𝑑u\mathrm{exp}\left[{\displaystyle _0^u}\left(A_vdW_v+(B_v{\displaystyle \frac{1}{2}}A_v^2)dv\right)\right][P_udW_u+(Q_uA_uP_u)du]\right),`$
which is the general solution of Eq. (40).

In Sec. VIII, we need only the case of Eq. (40) in which $`A_t=A`$, $`B_t=B`$, $`P_t=Pf_t`$, and $`Q_t=Qf_t`$, with $`A,B,P,Q`$ constants and with $`f_t`$ of the form $`f_t=\mathrm{exp}(Kt)`$, and so the solution of Eq. (A) then becomes

$`C_t`$ $`=`$ $`\mathrm{exp}\left(AW_t+(B{\displaystyle \frac{1}{2}}A^2)t\right)`$ (48)
$`\times `$ $`\left(C_0+{\displaystyle _0^t}𝑑u\mathrm{exp}\left(AW_u+(KB+{\displaystyle \frac{1}{2}}A^2)u\right)[PdW_u+(QAP)du]\right).`$
Using the identity (proved by the same methods used to find the integrating factor $`F_t`$),

$`\mathrm{exp}\left(\alpha W_u+(\beta +K)u\right)[PdW_u+(Q+\alpha P)du]`$
$`=`$ $`{\displaystyle \frac{P}{\alpha }}d\mathrm{exp}\left(\alpha W_u+(\beta +K)u\right){\displaystyle \frac{P}{\alpha }}(\beta +K{\displaystyle \frac{Q}{P}}\alpha {\displaystyle \frac{1}{2}}\alpha ^2)\mathrm{exp}\left(\alpha W_u+(\beta +K)u\right)du,`$
and taking $`\alpha =A`$ and $`\beta =B+A^2/2`$, the $`PdW_u`$ term in Eq. (48) can be eliminated. This gives an alternate form for the solution $`C_t`$,

$`C_t`$ $`=`$ $`\mathrm{exp}\left(AW_t+(B{\displaystyle \frac{1}{2}}A^2)t\right)`$ (50)
$`\times `$ $`(C_0{\displaystyle \frac{P}{A}}[\mathrm{exp}(AW_t+(KB+{\displaystyle \frac{1}{2}}A^2)t)1]`$
$`+`$ $`[{\displaystyle \frac{P}{A}}(KB)+Q]{\displaystyle _0^t}du\mathrm{exp}(AW_u+(KB+{\displaystyle \frac{1}{2}}A^2)u)).`$
Taking

$`A`$ $`=`$ $`\alpha _m^{(1)}={\displaystyle \frac{1}{2}}\sigma (E_mE_s),`$
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\alpha _m^{(1)})^2={\displaystyle \frac{1}{8}}\sigma ^2(E_mE_s)^2,`$
$`P`$ $`=`$ $`{\displaystyle \frac{1}{2}}\sigma V_{ms},`$
$`Q`$ $`=`$ $`iV_{ms}f_m,f_m=1{\displaystyle \frac{i}{8}}\sigma ^2(E_mE_s),`$
$`K`$ $`=`$ $`i(E_mE_sM){\displaystyle \frac{1}{2}}\mathrm{\Gamma },`$ (51)
in Eqs. (48) and (50) gives the results quoted respectively in Eqs. (25b) and (25) of the text.


format: text

text:
# Exchanging the places 𝑝 and ∞ in the Leopoldt conjecture

## 1 Introduction

It is often reasonable to treat the finite and the infinite places of a number field $`k`$ on an equal footing.
The unit group $`E=E_{\mathrm{}}(k)`$ of $`k`$ consists of all elements $`x`$ in $`k^{}`$ with normalized absolute value $`x_v=1`$ for every place $`v\mathrm{}`$. Fixing a prime $`p\mathrm{}`$ of $``$ we can equally well consider the group $`E_p(k)`$ of all $`x`$ in $`k^{}`$ with $`x_v=1`$ for all places $`vp`$. The elements of $`E_p(k)`$ are rational $`p`$-Weil numbers and it is known that $`E_p(k)`$ is finitely generated. The rank of $`E_p(k)`$ is known and also a way to construct a subgroup of finite index.

The easier part of the Dirichlet unit theorem asserts that the $`(\mathrm{log}_v)`$-map

(1)
$$E_{\mathrm{}}(k)H_{\mathrm{}}=\underset{v|\mathrm{}}{}v$$

induces an $``$-linear injection:

(2)
$$E_{\mathrm{}}(k)_{}H_{\mathrm{}}.$$

For $`E_p(k)`$ there is a corresponding map constructed by Gross \[G\]:

(3)
$$E_p(k)H_p=\underset{v|p}{}_pv.$$

In this case the injectivity of the analogue of (2)

(4)
$$E_p(k)_{}_pH_p$$

is only a conjecture. Gross proved it under certain conditions using Brumer’s $`p`$-adic analogue of Baker’s result on linear independence of logarithms of algebraic numbers. See \[G\] Cor. 2.14 or section two below.

Leopoldt’s conjecture is concerned with the $`p`$-adic logarithm

(5)
$$E_{\mathrm{}}(k)k_p=k_p.$$

Again the assertion is that its $`_p`$-linear extension

(6)
$$E_{\mathrm{}}(k)_{}_pk_p$$

is injective. Again using Brumer’s result this can be shown for abelian extensions $`k/`$.
In this note we will first review the preceeding facts in more detail. Then we consider the multivalued $`\mathrm{log}`$-map

(7)
$$E_p(k)k_{\mathrm{}}=k$$

which is a counterpart to the Leopoldt map (5) with the places $`p`$ and $`\mathrm{}`$ exchanged.
As it turns out the analogue of the Leopoldt conjecture predicts that certain vectors should be $``$-linear independent whose components are arguments of conjugates of Weil numbers in $`k`$. We give the precise formulation in 3.1. It is more involved than for the maps (2), (4) and (6) since the complex logarithm is multivalued. Using Baker’s result on linear forms in logarithms we prove part of our conjecture in certain abelian situations.

The multivalued map (7) gives rise to a single-valued map of $`E_p(k)`$ into a real torus. In section 4 we describe the closure of its image.

The Dirichlet regulator constructed via (2) is related to the zeta function of $`k`$ both at $`s=0`$ and $`s=1`$. The Gross regulator based on (4) is related to the $`p`$-adic zeta function of $`k`$ at $`s=0`$. The Leopoldt regulator based on (6) is related to its behaviour at $`s=1`$. It is an open question whether the regulators one can construct using (7) are connected to some kind of zeta function.

The relation of Weil numbers with motives over finite fields is explained in \[M\]. In \[GO\], \[O\] and \[K\], Weil numbers are studied from the point of view of algebraic and analytic number theory.

It is a pleasure for me to thank Niko Naumann and Kay Wingberg for helpful comments and the CRM in Montreal for its support.

## 2 Review of some material from algebraic number theory

For a number field $`k/`$ and a prime number $`p`$ we review some known facts about the group $`E_p(k)`$ defined in the introduction. Our basic reference is \[M\].
An algebraic number $`\lambda `$ is a $`q=p^r`$-Weil number of weight $`w`$ if it has the following two properties:

i) For every embedding $`\sigma :[\lambda ]`$ we have $`|\sigma (\lambda )|=q^{w/2}`$.
ii) There is some $`n`$ such that $`p^n\lambda `$ is an algebraic integer.

The set $`W_w(q)`$ of Weil numbers of weight $`w`$ in an algebraic closure $`\overline{}`$ of $``$ is stable under the action of $`G_{}=\mathrm{Gal}(\overline{}/)`$. The field generated by a set of $`q`$-Weil numbers is either totally real or a $`CM`$ field \[GO\] or \[M\] Def. 2.5 ff.

If $`k`$ contains a $`CM`$ field, let $`k_{CM}`$ denote the maximal $`CM`$ field contained in $`k`$. The following is clear:

###### Proposition 2.1

1) For any embedding $`k\overline{}`$ we have:

$$E_p(k)=W_0(p)k^{}.$$

2) If $`k`$ contains a $`CM`$ field, then

$$E_p(k)=E_p(k_{CM}).$$

If not then $`E_p(k)=\{\pm 1\}`$.

In the study of the unit group $`E=E_{\mathrm{}}(k)`$ one considers the regulator map obtained by composition

(8)
$$\alpha _{\mathrm{}}:E_{\mathrm{}}(k)\stackrel{\Delta }{}\underset{v|\mathrm{}}{}k_v^{}\stackrel{()_{v|\mathrm{}}}{}\underset{v|\mathrm{}}{}k_v^{}_v\stackrel{\underset{v}{}\mathrm{log}}{}\underset{v|\mathrm{}}{}v.$$

Here $`\mathrm{\Delta }`$ is the diagonal inclusion, $`k_v`$ is the $`v`$-completion and $`k_v^{}_v`$ is the value group $`_+^{}`$ of the absolute value $`_v`$. An analogous map is used for the study of $`E_p(k)`$ in \[M\] Prop. 2.27. It is the composition

(9)
$$\alpha _{(p)}:E_p(k)\stackrel{\Delta }{}\underset{v|p}{}k_v^{}\stackrel{(_v)_{v|p}}{}\underset{v|p}{}k_v^{}_v\stackrel{\underset{v}{}\mathrm{log}_{(p)}}{}\underset{v|p}{}v.$$

Here $`\mathrm{log}_{(p)}(r)=\mathrm{log}r/\mathrm{log}p`$ for $`rk_v^{}_v=Nv^{}`$. Thus $`\alpha _{(p)}`$ is given by

(10)
$$\alpha _{(p)}(x)=\underset{v|p}{}f(v|p)\mathrm{ord}_v(x)v.$$

By 2.1 we may assume that $`k`$ is a $`CM`$ field. We then denote the places dividing $`p`$ by $`𝔭`$ instead of $`v`$. Let $`c`$ denote complex conjugation on $`k`$. Then we have for $`xE_p(k)`$ that

$$\mathrm{ord}_𝔭(x)+\mathrm{ord}_{𝔭^c}(x)=\mathrm{ord}_𝔭(x)+\mathrm{ord}_𝔭(x^c)=\mathrm{ord}_𝔭(xx^c)=\mathrm{ord}_𝔭(1)=0.$$

Thus $`\alpha _{(p)}`$ takes values in

$$\left(\underset{𝔭|p}{}𝔭\right)^{}=\left\{\underset{𝔭|p}{}\nu _𝔭𝔭\right|\nu _𝔭+\nu _{𝔭^c}=0\text{for all}𝔭\}.$$

Besides the map

(11)
$$\alpha _{(p)}:E_p(k)\left(\underset{𝔭|p}{}𝔭\right)^{}$$

there are maps in the other direction. Let $`𝒯`$ be the set of primes $`𝔭|p`$ in $`k`$ with $`𝔭𝔭^c`$. Choose an integer $`M1`$ which is divisible by all inertia indeces $`f(𝔭|p)`$ for $`𝔭`$ in $`𝒯`$ and such that $`𝔭^{Mf(𝔭|p)^1}`$ is a principal ideal for these $`𝔭`$. We may then choose $`x_𝔭`$ in $`𝔬_k`$ for $`𝔭`$ in $`𝒯`$ such that we have $`𝔭^{Mf(𝔭|p)^1}=(x_𝔭)`$ and $`x_𝔭^c=x_{𝔭^c}`$. Define a map

(12)
$$\pi _M:\left(\underset{𝔭|p}{}𝔭\right)^{}E_p(k)$$

by the formula

$$\pi _M\left(\underset{𝔭|p}{}\nu _𝔭𝔭\right)=\underset{𝔭𝒯}{}x_𝔭^{\nu _𝔭}.$$

Essentially, the following facts can be found in \[S\] Ch. II or in \[M\]. For convenience we give the short proof.

###### Proposition 2.2

Let $`k`$ be a $`CM`$ field and let $`S`$ be a system of representatives of the primes in $`𝒯`$ under complex conjugation. The following assertions hold:
1) $`\alpha _{(p)}\pi _M=M\mathrm{id}`$ and $`(\pi _M\alpha _{(p)})(x)x^Mmod\mu (k)`$.
2) The abelian group $`E_p(k)`$ is finitely generated with

$`\mathrm{tors}E_p(k)`$ $`=`$ $`\mu (k)=\mathrm{Ker}\alpha _{(p)}\text{and}`$
$`\mathrm{rank}E_p(k)`$ $`=`$ $`{\displaystyle \frac{1}{2}}|𝒯|{\displaystyle \frac{1}{2}}(k:).`$
3) Rationally $`\alpha _{(p)}`$ and $`\pi _M`$ are isomorphisms.
4) The elements $`\xi _𝔭=x_𝔭^cx_𝔭^1E_p(k)`$ for $`𝔭S`$ form a $``$-basis of $`E_p(k)`$.

Proof 1) The first relation follows from a straightforward computation. As for the second, set

$$y=x^M\pi _M(\alpha _{(p)}(x))=x^M\underset{𝔭𝒯}{}x_𝔭^{f(𝔭|p)\mathrm{ord}_𝔭(x)}.$$

Then $`(y)=(x)^M_{𝔭|p}𝔭^{M\mathrm{ord}_𝔭(x)}=(1)`$. Hence $`y`$ is a unit.

Since $`y`$ lies in $`E_p(k)`$ it follows that all absolute values of $`y`$ are equal to $`1`$. Thus $`y`$ is a root of unity.
2, 3) Since $`\alpha _{(p)}`$ embeds $`E_p(k)/\mu (k)`$ injectively into a finitely generated abelian group, $`E_p(k)`$ is itself finitely generated. By 1) the map $`\alpha _{(p)}`$ is an isomorphism rationally, hence

$$\mathrm{rank}E_p(k)=dim\left(\underset{𝔭|p}{}𝔭\right)^{}=\frac{1}{2}|𝒯|.$$

4) The group $`\left(_{𝔭|p}𝔭\right)^{}`$ has basis the elements $`𝔭^c𝔭`$ for $`𝔭S`$. They are mapped to $`x_{𝔭^c}x_𝔭^1=x_𝔭^cx_𝔭^1`$ under $`\pi _M`$. Now 3) implies the assertion.  $`\mathrm{}`$

Apart from (9) there is a deeper $`p`$-analogue of the Dirichlet regulator map $`\alpha _{\mathrm{}}`$ in (8). It was introduced by Gross in order to formulate a $`p`$-adic analogue of the Stark conjectures \[G\].

Gross first defines “absolute values” with values in $`_p`$” for every valuation $`v`$ of $`k`$. Namely, for $`x`$ in $`k_v^{}`$ set

$$x_{v,p}=\{\begin{array}{cc}1\hfill & v\text{complex}\hfill \\ \mathrm{sign}x\hfill & v\text{real}\hfill \\ (Nv)^{\mathrm{ord}_v(x)}\hfill & v\text{finite}vp\hfill \\ (Nv)^{\mathrm{ord}_v(x)}N_{k_v/_p}(x)\hfill & v|p.\hfill \end{array}$$

As for the usual absolute values the product formula holds:

$$\underset{v}{}x_{v,p}=1\text{for}xk^{}.$$

Let $`\mathrm{log}_p:_p^{}_p`$ be the $`p`$-adic logarithm. Then Gross’ $`p`$-adic analogue of $`\alpha _{\mathrm{}}`$ is the composition:

(13)
$$\alpha _p:E_p(k)\stackrel{\Delta }{}\underset{v|p}{}k_v^{}\stackrel{(_{v,p})_{v|p}}{}\underset{v|p}{}k_v^{}_{v,p}\stackrel{\underset{v}{}\mathrm{log}_p}{}\underset{v|p}{}_pv.$$

Its kernel consists of the roots of unity in $`k`$ and Gross conjectures that tensored with $`_p`$ the map $`\alpha _p`$ remains injective, \[G\] conjecture 1.15.

In \[G\] corollary 2.14 this is proved in the following case: $`k`$ contains a $`CM`$ field and the set of places $`𝔭^+|p`$ in $`k_{CM}^+`$ which split in $`k_{CM}`$ is permuted transitively by an abelian subgroup of $`\mathrm{Aut}k_{CM}^+`$.

We now turn to the Leopoldt conjecture whose analogue for $`p`$ and $`\mathrm{}`$ exchanged will be studied in the next section.

For every place $`v|p`$ of $`k`$ let $`\mathrm{log}_v:k_v^{}k_v`$ be the $`p`$-adic logarithm normalized by $`\mathrm{log}_vp=1`$. Consider the composition:

(14)
$$\alpha _{\mathrm{},p}:E_{\mathrm{}}(k)\stackrel{\Delta }{}\underset{v|p}{}k_v^{}\stackrel{\underset{v}{}\mathrm{log}_v}{}\underset{v|p}{}k_v=k_p.$$

Its kernel consists of roots of unity, so that

$$\alpha _{\mathrm{},p}1:E_{\mathrm{}}(k)k_p$$

is injective. The Leopoldt conjecture asserts that the map

$$\alpha _{\mathrm{},p}1:E_{\mathrm{}}(k)_pk_p$$

is injective. It is known for abelian extensions $`k/`$ for example \[W\].

## Appendix

Let $`\lambda `$ be a $`q`$-Weil number of weight $`w`$ in $`\overline{}`$. Let $`k`$ be a $`CM`$ field containing $`\lambda `$. Then $`\lambda `$ is determined up to a root of unity by the valuations $`\lambda _𝔭`$ for $`𝔭|p`$ in $`k`$.

If $`\alpha `$ is such that $`q^\alpha =\lambda `$, then $`\mathrm{Re}\alpha =\frac{w}{2}`$. The imaginary part of $`\alpha `$ is determined by $`\lambda `$ up to an element of $`\frac{2\pi }{\mathrm{log}q}`$. Hence the values $`\lambda _𝔭`$ for $`𝔭|p`$ determine $`\mathrm{Im}\alpha mod\frac{2\pi }{\mathrm{log}q}`$. The material of the first section allows one to write down a formula for this dependence. For a variety $`X/𝔽_q`$ this gives a way to reflect $`𝔭`$-adic divisibility properies of the zeroes of the zeta function $`Z_X(T)`$ in the imaginary parts of the zeroes of the Hasse Weil zeta function $`\zeta _X(s)`$. Recall here that

$$\zeta _X(s)=\underset{x|X|}{}(1Nx^s)^1\text{for}s>dimX\text{in}$$

and that the two zeta functions are related by the formula:

$$\zeta _X(s)=Z_X(q^s).$$

In particular, I wonder whether a statement such as “the Newton polygon lies above the Hodge polygon” might be reformulated in terms of $`\zeta _X(s)`$ and its zeroes such that it makes sense for $`X/`$.

However that may be, the desired formula for $`\mathrm{Im}\alpha `$ in $`/\frac{2\pi }{\mathrm{log}q}`$ is the following:

$`\stackrel{~}{\mathrm{Im}}\alpha `$ $`=`$ $`{\displaystyle \underset{𝔭𝒯}{}}\mathrm{log}_{(p)}\lambda _𝔭\stackrel{~}{\mathrm{arg}}_q(x_𝔭^{1/M})`$
$`=`$ $`{\displaystyle \underset{𝔭𝒯}{}}f(𝔭|p)\mathrm{ord}_𝔭(\lambda )\stackrel{~}{\mathrm{arg}}_q(x_𝔭^{1/M}).`$
Here $`\mathrm{arg}_q=(\mathrm{log}q)^1\mathrm{arg}`$ and for a real number $`t`$ we let $`\stackrel{~}{t}`$ be its image in $`/\frac{2\pi }{\mathrm{log}q}`$. The integer $`M`$ is chosen as in section 2 before proposition 2.2.

Proof The number $`x=\lambda /\lambda ^c`$ is in $`E_p(k)`$. By proposition 2.2 we have

(16) $`2\stackrel{~}{\mathrm{arg}}_q\lambda `$ $`=`$ $`\stackrel{~}{\mathrm{arg}}_qx=M^1\stackrel{~}{\mathrm{arg}}(\pi _M\alpha _{(p)}(x))`$
$`=`$ $`M^1\stackrel{~}{\mathrm{arg}}_q\pi _M\left({\displaystyle \underset{𝔭𝒯}{}}\mathrm{log}_{(p)}\lambda /\lambda ^c_𝔭𝔭\right)`$
$`=`$ $`M^1{\displaystyle \underset{𝔭𝒯}{}}\mathrm{log}_{(p)}\lambda /\lambda ^c_𝔭\stackrel{~}{\mathrm{arg}}_qx_𝔭.`$
Since $`x_{\overline{𝔭}}=\overline{x}_𝔭`$ we find

$$\stackrel{~}{\mathrm{arg}}_q\lambda =M^1\underset{𝔭𝒯}{}\mathrm{log}_{(p)}\lambda _𝔭\stackrel{~}{\mathrm{arg}}_qx_𝔭.$$

Hence (Appendix) follows.  $`\mathrm{}`$

Another way to write $`\stackrel{~}{\mathrm{Im}}\alpha `$ is the following: Let $`S`$ be as in proposition 2.2. Then starting from (16) we find

$$\stackrel{~}{\mathrm{Im}}\alpha =\underset{𝔭S}{}\mathrm{log}_{(p)}(\lambda _𝔭/\lambda _{𝔭^c})\stackrel{~}{\mathrm{arg}}_q(\xi _𝔭)\text{in}/\frac{2\pi }{\mathrm{log}q}$$

where $`\xi _𝔭=(x_𝔭^cx_𝔭^1)^{1/M}`$ in $`^{}/\mu ()`$.

In the $``$-vector space $`/\frac{2\pi }{\mathrm{log}q}`$ the vectors $`\stackrel{~}{\mathrm{arg}}_q(\xi _𝔭)`$ for $`𝔭`$ in $`S`$ are $``$-linearly independent by Prop. 2.2, 4). Hence we see how $`\stackrel{~}{\mathrm{Im}}\alpha `$ determines the quotients $`\lambda _𝔭/\lambda _{𝔭^c}`$ for all $`𝔭S`$. These in turn determine $`\lambda _𝔭`$ for all $`𝔭|p`$ since $`\lambda `$ was a $`q`$-Weil number of weight $`w`$.

## 3 Changing places in the Leopoldt conjecture

In this section we deal with an analogue of the Leopoldt regulator map (14) for $`E_p(k)`$ instead of $`E_{\mathrm{}}(k)`$. We may therefore assume that $`k`$ is a $`CM`$ field. Ideally we would consider the composition:

(17)
$$E_p(k)\stackrel{\Delta }{}\underset{v|\mathrm{}}{}k_v^{}\stackrel{\underset{v}{}\mathrm{log}_v}{}\underset{v|\mathrm{}}{}k_v=k.$$

Here $`\mathrm{log}_v`$ is the complex logarithm on $`k_v^{}^{}`$. Unfortunately the complex logarithm is a multivalued map, so (17) has to be interpreted suitably.

Note that for $`x`$ in $`E_p(k)`$ we have

$$2\mathrm{R}\mathrm{e}\mathrm{log}_vx=\mathrm{log}x_v=\mathrm{log}1=0.$$

We set

$$(1)_v=\{zk_v|\mathrm{Re}z=0\}.$$

More generally, for any subgroup $`\mathrm{\Lambda }`$ let us write $`\mathrm{\Lambda }(1)_v`$ for the image of $`2\pi i\mathrm{\Lambda }`$ in $`k_v`$ under any of the two continuous isomorphisms $`\stackrel{}{}k_v`$. Then we interpret (17) as the injective map:

(18)
$$\alpha _{p,\mathrm{}}:E_p(k)\stackrel{\Delta }{}\underset{v|\mathrm{}}{}k_{v,1}^{}\stackrel{\underset{v}{}\mathrm{log}_v}{}T_k=\underset{v|\mathrm{}}{}(1)_v/(1)_v.$$

Here $`k_{v,1}^{}`$ is the subgroup of $`k_v^{}`$ of elements with norm one.

Since $`T_k`$ is not an $``$-vector space we cannot state a “Leopoldt conjecture” in quite the same way as before. The strongest formulation seems to be the following.

###### Conjecture 3.1

Let $`\xi _1,\mathrm{},\xi _m`$ be a basis of $`E_p(k)`$ and consider the diagram

$$E_p(k)\stackrel{\alpha _{p,\mathrm{}}1}{}\underset{v|\mathrm{}}{}(1)_v/(1)_v\stackrel{pr}{}V_k=\underset{v|\mathrm{}}{}(1)_v.$$

For any $`1im`$ choose a preimage $`\eta _i`$ under $`\mathrm{pr}`$ of $`(\alpha _{p,\mathrm{}}1)(\xi _i)`$. Then the vectors $`\eta _i`$ are $``$-linearly independent in $`V_k`$.

Note here that

$$\mathrm{rank}E_p(k)=|S|\frac{1}{2}(k:)=dimV_k.$$

We have equality if and only if $`p`$ splits completely in $`k`$. In this case one may form regulator determinants. Observe also that we get a weaker conjecture if we work integrally instead of tensoring with $``$.

###### Proposition 3.2

It is sufficient to verify the conjecture for any fixed basis $`\xi _1,\mathrm{},\xi _m`$ (but all possible liftings $`\eta _i`$).

The proof is straightforward.

Choose a representative $`\sigma _v:k`$ for every complex place $`v`$ of $`k`$. Let $`S`$ be as in proposition 2.2. Then according to propositions 2.2, 4) and 3.2 our conjecture is equivalent to the following:

###### Conjecture 3.3

The vectors $`(\mathrm{arg}(\overline{\sigma _v(x_𝔭)}\sigma _v(x_𝔭)^1))_{v|\mathrm{}}`$ for $`𝔭S`$ are $``$-linearly independent in $`^{r_2}`$ for any choices of argument $`mod2\pi `$.

Trivial example Assume that there are exactly two primes dividing $`p`$ in $`k`$. They are $`𝔭`$ and $`𝔭^c𝔭`$ for some $`𝔭|p`$. Take $`S=\{𝔭\}`$. The conjecture asserts that for some $`v|\mathrm{}`$

$$\mathrm{arg}(\overline{\sigma _v(x_𝔭)}\sigma _v(x_𝔭)^1)2\pi .$$

In other words $`\overline{\sigma _v(x_𝔭)}\sigma _v(x_𝔭)^1`$ should not be a root of unity. This is clear, since otherwise we would have $`(𝔭^c)^N=𝔭^N`$ for some $`N1`$, a contradiction.

Next we introduce special bases of $`E_p(k)`$, at least in favourable circumstances.

###### Proposition 3.4

Let $`k`$ be a $`CM`$ field and consider as before the set $`𝒯`$ of primes $`𝔭`$ dividing $`p`$ in $`k`$ such that $`𝔭^c𝔭`$. Assume that $`\mathrm{Aut}k`$ contains a subgroup $`H`$ with $`cH`$ such that $`G=H,c`$ permutes the primes in $`𝒯`$ transitively with isotropy subgroups contained in $`H`$.
Then there exists an element $`\xi `$ in $`E_p(k)`$ such that the set of conjugates $`\xi ^\sigma `$ for $`\sigma `$ in $`H`$ is a basis of $`E_p(k)`$.

Remarks 1 The conjugates $`\xi ^\sigma `$ do not need to be pairwise different.
2 As $`c`$ is central we have $`GH\times c`$ and $`H`$ is normal in $`G`$. Hence the condition on isotropy subgroups has to be checked for one prime in $`𝒯`$ only. It is trivially satisfied if $`p`$ is completely decomposed in $`k`$. In that case $`k/`$ must be Galois and the conjugates of $`\xi `$ under $`H`$ are pairwise distinct.

Proof We may assume that $`𝒯`$ is non-empty. Fix a prime $`𝔭_0`$ in $`𝒯`$, let $`k^{}`$ be the fixed field of $`G`$ and set $`𝔭_0^{}=𝔭_0k^{}`$. Then we claim:

(19)
$$𝒯=\text{set of primes}𝔭\text{in}k\text{dividing}𝔭_0^{}.$$

Indeed, if $`𝔭𝔭^c`$ divides $`p`$, then by assumption on $`G`$ there is some $`\sigma G`$ with $`𝔭=𝔭_0^\sigma `$. Hence $`𝔭k^{}=(𝔭_0k^{})^\sigma =𝔭_0^{}`$. On the other hand, the prime divisors $`𝔭`$ of $`𝔭_0^{}`$ are all conjugate under $`G`$ since the extension $`k/k^{}`$ is Galois with group $`G`$. For any $`𝔭|𝔭_0^{}`$ there is thus some $`\sigma `$ in $`G`$ with $`𝔭=𝔭_0^\sigma `$. Hence $`𝔭^c=𝔭_0^{\sigma c}=𝔭_0^{c\sigma }𝔭_0^\sigma =𝔭`$ since $`𝔭_0^c𝔭_0`$.

By proposition 2.2 the map $`\alpha _{(p)}`$ in (11) induces an isomorphism

(20)
$$\alpha _{(p)}:E_p(k)\stackrel{}{}\left(\underset{𝔭|p}{}𝔭\right)^{}=[𝒯]^{}$$

where $`()^{}`$ denotes the $`1`$ eigenspace of the action of complex conjugation. The map $`\alpha _{(p)}`$ is $`G`$-equivariant under the natural left $`G`$-actions on both sides.

Now let $`G_{𝔭_0}`$ be the decomposition group of $`𝔭_0`$ in the extension $`k/k^{}`$. According to (19) we have a bijection

(21)
$$G/G_{𝔭_0}\stackrel{}{}𝒯,\sigma G_{𝔭_0}\sigma (𝔭_0).$$

Together with (20) we obtain a $`G`$-equivariant isomorphism:

(22)
$$E_p(k)\stackrel{}{}[G/G_{𝔭_0}]^{}.$$

By assumption $`G_{𝔭_0}H`$. Hence we get an embedding

$$H/G_{𝔭_0}G/G_{𝔭_0}c\times H/G_{𝔭_0}.$$

The natural restriction map

$$\mathrm{res}:[G/G_{𝔭_0}][H/G_{𝔭_0}]$$

induces an $`H`$-equivariant isomorphism

$$\mathrm{res}:[G/G_{𝔭_0}]^{}\stackrel{}{}[H/G_{𝔭_0}].$$

Combined with (22) we get an $`H`$-equivariant isomorphism:

(23)
$$E_p(k)\stackrel{}{}[H/G_{𝔭_0}].$$

Hence any element $`\xi `$ in $`E_p(k)`$ corresponding to a multiple of $`\delta _e=1eG_{𝔭_0}`$ has the required property.

More explicitely we can describe such a $`\xi `$ as follows. Let $`k_{𝔭_0}`$ be the fixed field of $`G_{𝔭_0}`$. In the extension $`k/k_{𝔭_0}`$ only one prime lies above $`𝔮_0=𝔭_0k_{𝔭_0}`$, namely $`𝔭_0`$. Thus we have $`𝔮_0𝔬_k=𝔭_0^e`$ where $`e1`$. Choose $`yk_{𝔭_0}`$ such that $`(y)=𝔮_0^h`$ in $`k_{𝔭_0}`$ for some $`h1`$. It follows that $`(y)=𝔭_0^{eh}`$ in $`k`$. Choose a multiple $`M1`$ of $`efh`$ where $`f`$ is the common inertia index over $`p`$ of the primes in $`𝒯`$ and set $`x=y^{H/efh}k_{𝔭_0}`$. Then we have $`(x)=𝔭_0^{Hf^1}`$ in $`k`$. The $`H/G_{𝔭_0}`$ conjugates of $`\xi =x^cx^1E_p(k_{𝔭_0})E_p(k)`$ give a basis of $`E_p(k)`$.  $`\mathrm{}`$

###### 3.5

We now prove special cases of conjecture 3.1. Let $`k`$ be a $`CM`$ field and assume that we are given an element $`\xi `$ in $`E_p(k)`$ and an abelian subgroup $`H`$ of $`\mathrm{Aut}k`$ with $`cH`$ such that the set of conjugates $`\xi ^\sigma `$ for $`\sigma H`$ is a basis of $`E_p(k)`$.

Let $`H_\xi `$ be the group of $`\sigma H`$ fixing $`\xi `$ and let $`\overline{H}`$ be the abelian quotient group $`\overline{H}=H/H_\xi `$. Fix an embedding $`k`$.

For every $`\overline{\sigma }\overline{H}`$ choose a preimage denoted by $`\sigma _v`$ in $`H`$ and hence in $`\mathrm{Aut}k`$. Via the embedding $`k`$ we will also view $`\sigma _v`$ as an embeding $`\sigma _v:k`$. Let $`v`$ be the infinite place corresponding to $`\sigma _v`$. Then $`\sigma _v`$ induces isomorphisms $`\sigma _v:k_v\stackrel{}{}`$ and $`\sigma _v:(1)_v\stackrel{}{}(1)`$. Let $`S_{\mathrm{}}`$ be the set of infinite places so obtained. The bijection $`\overline{H}\stackrel{}{}\{\xi ^\sigma |\sigma H\},\overline{\tau }\overline{\tau }(\xi )`$ shows that

$$|S_{\mathrm{}}|=|\overline{H}|=\mathrm{rank}E_p(k).$$

Next, let $`s`$ be a set theoretical splitting of the projection $`(1)(1)/(1)`$. It induces splittings $`s_v`$ of $`(1)_v(1)_v/(1)_v`$ via conjugation, $`s_v=\sigma _v^1s\sigma _v`$.

For each basis vector $`\overline{\tau }(\xi )`$ of $`E_p(k)`$ with $`\overline{\tau }\overline{H}`$ define as follows a lift $`\eta _{\overline{\tau }}`$ to $`V_k`$ of the vector

$$(\alpha _{p,\mathrm{}}1)(\overline{\tau }(\xi ))=(\mathrm{log}_v(\overline{\tau }(\xi )_v))_{v|\mathrm{}}.$$

The $`v`$-components with $`v`$ in $`S_{\mathrm{}}`$ are lifted via $`s_v`$. The other components are lifted arbitrarily. For $`v`$ in $`S_{\mathrm{}}`$ we therefore have:

$$\eta _{\overline{\tau },v}=s_v\mathrm{log}_v(\overline{\tau }(\xi )_v).$$

###### Theorem 3.6

In the situation of 3.5 the vectors $`\eta _{\overline{\tau }}`$ for $`\overline{\tau }\overline{H}`$ are $``$-linearly independent in $`V_k`$.

Proof It suffices to show that

$$\delta =|det(\sigma _vs_v\mathrm{log}_v(\overline{\tau }(\xi )_v))_{\overline{\tau }\overline{H},vS_{\mathrm{}}}|$$

is non-zero. Note that the determinant itself is well defined only up to sign.

We have

$`\sigma _vs_v\mathrm{log}_v(\overline{\tau }(\xi )_v)`$ $`=`$ $`s\sigma _v\mathrm{log}_v(\overline{\tau }(\xi )_v)=s\mathrm{log}\sigma _v\overline{\tau }(\xi )`$
$`=`$ $`s\mathrm{log}\overline{\sigma }_v\overline{\tau }(\xi ).`$
We may order $`S_{\mathrm{}}`$ and $`\overline{H}`$ in such a way that we get:

$$\delta =|det(s\mathrm{log}\overline{\sigma }^1\overline{\tau }(\xi ))_{\overline{\sigma },\overline{\tau }\overline{H}}|.$$

The well known formula for the group determinant \[W\] Lemma 5.26 gives:

$$\delta =\underset{\chi \overline{H}^{}}{}\left|\underset{\overline{\sigma }\overline{H}}{}\chi (\overline{\sigma })s\mathrm{log}\overline{\sigma }(\xi )\right|.$$

We may choose $`N1`$ so that

$$Ns\mathrm{log}\overline{\sigma }(\xi )=\mathrm{log}_{\overline{\sigma }}(\overline{\sigma }(\xi )^N)\text{for all}\overline{\sigma }\overline{H}.$$

Here $`\mathrm{log}_{\overline{\sigma }}`$ is a suitable branch of the complex logarithm depending on $`\overline{\sigma }`$. It suffices therefore to show that

$$\underset{\overline{\sigma }\overline{H}}{}\chi (\overline{\sigma })\mathrm{log}_{\overline{\sigma }}(\overline{\sigma }(\xi )^N)0\text{for every character}\chi .$$

This is an immediate consequence of Baker’s theorem \[B\] Theorem 2.1, since the numbers $`\mathrm{log}_{\overline{\sigma }}(\overline{\sigma }(\xi )^N)`$ are $``$-linear independent as one sees upon exponentiating.  $`\mathrm{}`$

###### Corollary 3.7

Let $`k`$ be a $`CM`$ field satisfying the assumptions of proposition 3.4 with $`H`$ abelian and choose any $`\xi `$ constructed there. We may then apply the construction of 3.5 to get liftings $`\eta _{\overline{\tau }}`$ of $`(\alpha _{p,\mathrm{}}1)(\overline{\tau }(\xi ))`$ for $`\overline{\tau }\overline{H}`$ which are linearly independent in $`V_k`$.

This lends some credibility to the conjecture. In general, even in the abelian case it seems beyond the power of Baker’s theorem. For example assume we are in the situation of proposition 3.4 with $`H=\{1,\sigma \}`$ of order two and exactly four different primes of $`k`$ dividing $`p`$. Then $`\xi ,\xi ^\sigma `$ form a basis of $`E_p(k)`$. Baker’s theorem allows us to show that for any two branches $`\mathrm{log}^{(1)}`$ and $`\mathrm{log}^{(2)}`$ of the complex logarithm, the determinant

$$\left|\begin{array}{cc}\mathrm{log}^{(1)}\xi & \mathrm{log}^{(2)}\xi ^\sigma \\ \mathrm{log}^{(2)}\xi ^\sigma & \mathrm{log}^{(1)}\xi \end{array}\right|$$

is non-zero. However if we choose the branches of the logarithm less coherently we are no longer reduced to linear forms in logarithms. For example, I cannot prove the following determinant to be non-zero:

$$\left|\begin{array}{cc}\mathrm{log}^{(1)}\xi & \mathrm{log}^{(2)}\xi ^\sigma \\ \mathrm{log}^{(2)}\xi ^\sigma & \mathrm{log}^{(3)}\xi \end{array}\right|$$

where $`\mathrm{log}^{(3)}=\mathrm{log}^{(1)}+2\pi i`$.

Let us finally discuss a case where $`k/`$ is abelian but the requirement in proposition 3.4 that $`c`$ has a complement is not met. Take $`k=(\zeta _5)`$. Its Galois group is cyclic of order $`4`$. Let $`\sigma _\nu `$ be the automorphism corresponding to $`\nu `$ in $`(/5)^{}`$.

We are interested in the prime $`p=11`$. It decomposes into the primes
$`𝔭_\nu =(1+2\zeta _5^\nu )`$ for $`\nu (/5)^{}`$. Thus $`𝒯=\{𝔭_\nu |\nu (/5)^{}\}`$. Let us choose $`S=\{𝔭_1,𝔭_2\}`$ for example. We may take $`M=1`$ in the construction before proposition 2.2 and

$$x_{𝔭_1}=1+2\zeta _5\text{and}x_{𝔭_2}=1+2\zeta _5^2.$$

Then

$$\xi _{𝔭_1}=x_{𝔭_1}^cx_{𝔭_1}^1\text{and}\xi _{𝔭_2}=x_{𝔭_2}^cx_{𝔭_2}^1=\xi _{𝔭_1}^{\sigma _2}$$

form a basis of $`E_{11}(k)`$. Now let us use the same branch $`\mathrm{log}`$ of logarithm for the liftings in conjecture 3.1. Since $`\sigma _2^2=\sigma _1`$ instead of $`\sigma _2^2=1`$ we do not end up with a group determinant. Instead we get a determinant like the following

$$\left|\begin{array}{cc}\mathrm{log}\xi _{𝔭_1}& \mathrm{log}\xi _{𝔭_1}^{\sigma _2}\\ \mathrm{log}\xi _{𝔭_1}^{\sigma _2}& \mathrm{log}\xi _{𝔭_1}^{\sigma _2^2}\end{array}\right|=\left|\begin{array}{cc}\mathrm{log}\xi _{𝔭_1}& \mathrm{log}\xi _{𝔭_2}\\ \mathrm{log}\xi _{𝔭_2}& \mathrm{log}\xi _{𝔭_1}^1\end{array}\right|.$$

For an unfortunate choice of $`\mathrm{log}`$ we have $`\mathrm{log}\xi _{𝔭_1}^1=\mathrm{log}\xi _{𝔭_1}\pm 2\pi i`$ so that the determinant does not factor into a product of linear forms in logarithms. We have therefore not persued the case where $`c`$ does not have a complement $`H`$.

## 4 On the closure of $`\alpha _{p,\mathrm{}}(E_p(k))`$ in $`T_k`$

The Leopoldt conjecture can be stated as an assertion about the closure of the global units embedded into the local units at $`p`$. Regarding the map $`\alpha _{p,\mathrm{}}`$ in (18) one may ask about the closure of the subgroup $`\alpha _{p,\mathrm{}}(E_p(k))`$ in the real $`\frac{1}{2}(k:)`$-dimensional torus $`T_k`$. We have

$$\overline{\alpha _{p,\mathrm{}}(E_p(k))}=\alpha _{p,\mathrm{}}(E_p(k)\mu (k))\times \overline{\alpha _{p,\mathrm{}}(E_p(k))}^0$$

where <sup>0</sup> denotes the connected component of zero. In order to determine $`\overline{\alpha _{p,\mathrm{}}(E_p(k))}^0`$ it is equivalent but more convenient to describe the quotient

$$C=T_k/\overline{\alpha _{p,\mathrm{}}(E_p(k))}^0.$$

For the Pontrjagin dual we have:

$`\widehat{C}`$ $`=`$ $`\mathrm{Ker}\left(\widehat{T}_k(\overline{\alpha _{p,\mathrm{}}(E_p(k))}^0)^{}\right)`$
$`=`$ $`\mathrm{Ker}\left(\widehat{T}_k\alpha _{p,\mathrm{}}(E_p(k))^{}\right).`$
For simplicity assume that $`k/`$ is a Galois $`CM`$ extension with group $`G`$. Fix an embedding $`k`$ and for every $`v|\mathrm{}`$ choose some $`\sigma _vG`$ which gives rise to the place $`v`$ after composing with the inclusion $`k`$.

There is an isomorphism

$$\underset{v|\mathrm{}}{}\stackrel{}{}\widehat{T}_k,\mu \chi _\mu $$

with $`\chi _\mu (t)=_{v|\mathrm{}}\mathrm{exp}\mu _v(\sigma _vt_v)`$.

Here $`\sigma _v`$ is viewed as the isomorphism $`\sigma _v:(1)_v/(1)_v(1)/(1)`$ induced by $`\sigma _v:k_v\stackrel{}{}`$.

Since

$$\sigma _v\mathrm{log}_v\xi _v=\mathrm{log}\sigma _v\xi \text{for all}\xi E_p(k)$$

we get

$`\widehat{C}`$ $`=`$ $`\left\{\mu \text{in}{\displaystyle \underset{v|\mathrm{}}{}}\right|{\displaystyle \underset{v|\mathrm{}}{}}(\sigma _v\xi )^{\mu _v}\mu ()\text{for all}\xi \text{in}E_p(k)\}`$
$`=`$ $`\left\{\mu \text{such that}{\displaystyle \underset{v|\mathrm{}}{}}\mu _v\sigma _v\text{annihilates}E_p(k)\right\}.`$
Using the isomorphism (20) we find

$$\widehat{C}=\left\{\mu \text{such that}\underset{v|\mathrm{}}{}\mu _v\sigma _v\text{annihilates}[𝒯]^{}\right\}.$$

Set $`\epsilon =\frac{1}{2}(1c)`$. Then $`[𝒯]^{}=\epsilon [𝒯]`$ and we see that $`\widehat{C}`$ consists of those $`\mu `$ such that $`_{v|\mathrm{}}\mu _v(\sigma _vc\sigma _v)`$ annihilates $`[𝒯]`$ i.e. each $`𝔭𝒯`$. For $`𝔭,𝔭^{}`$ in $`𝒯`$ define the function

$$\epsilon _{𝔭,𝔭^{}}:\{v|\mathrm{}\}\{1,1,0\}$$

by

$$\epsilon _{𝔭,𝔭^{}}(v)=\{\begin{array}{cc}1\hfill & \text{if}\sigma _v𝔭=𝔭^{}\hfill \\ 1\hfill & \text{if}\sigma _v𝔭=c𝔭^{}\hfill \\ 0\hfill & \text{otherwise}.\hfill \end{array}$$

Define an inner product on $`_{v|\mathrm{}}`$ by the formula

$$a_vv,b_vv=a_vb_v.$$

Now an easy calculation gives the first part of the following theorem if we take into account the relations

$$\epsilon _{𝔭,𝔭^{}}=\epsilon _{c𝔭,c𝔭^{}}=\epsilon _{𝔭,c𝔭^{}}=\epsilon _{c𝔭,𝔭^{}}.$$

We let $`S𝒯`$ be as in proposition 2.2.

###### Theorem 4.1

Let $`k/`$ be a Galois $`CM`$ extension. Then we have:

$$\widehat{C}=\left\{\mu \text{in}\underset{v|\mathrm{}}{}\text{such that}\mu ,\epsilon _{𝔭,𝔭^{}}=0\text{for all}𝔭,𝔭^{}S\right\}.$$

Moreover

$$dim\overline{\alpha _{p,\mathrm{}}(E_p(k))}=\begin{array}{c}\text{dimension of the}\text{-span of the elements}\epsilon _{𝔭,𝔭^{}}\hfill \\ \text{for}𝔭,𝔭^{}S\text{in}_{v|\mathrm{}}.\hfill \end{array}$$

In particular $`\alpha _{p,\mathrm{}}(E_p(k))`$ is dense in $`T_k`$ if $`p`$ is completely decomposed.

Proof We have

$`dim\overline{\alpha _{p,\mathrm{}}(E_p(k))}`$ $`=`$ $`dim\overline{\alpha _{p,\mathrm{}}(E_p(k))}^0=dimT_kdimC`$
$`=`$ $`dimT_kdim_{}(\widehat{C})`$
$`=`$ $`dimT_kdim_{}\epsilon _{𝔭,𝔭^{}},𝔭,𝔭^{}\text{in}S^{}`$
$`=`$ $`dim_{}\epsilon _{𝔭,𝔭^{}},𝔭,𝔭^{}\text{in}S.`$
If $`p`$ is completely decomposed, then $`\sigma _v𝔭=\sigma _w𝔭`$ implies $`v=w`$. Moreover $`\sigma _v𝔭=c\sigma _w𝔭`$ is not possible for any two places $`v,w`$. Hence

$$\epsilon _{𝔭,𝔭^{\sigma _w}}(v)=\{\begin{array}{ccc}1\hfill & \text{if}\hfill & v=w\hfill \\ 0\hfill & \text{if}\hfill & vw.\hfill \end{array}$$

Hence the elements $`\epsilon _{𝔭,𝔭^{\sigma _w}}\mathrm{span}_{v|\mathrm{}}`$ and therefore $`\overline{\epsilon _{p,\mathrm{}}(E_p(k))}`$ is all of $`T_k`$.  $`\mathrm{}`$

Mathematisches Institut
Westf. Wilhelms-Universität
Einsteinstr. 62
48149 Münster
Germany
deninge@math.uni-muenster.de


text:
# REFERENCES

Non-factorized genuine twist $`3`$ in exclusive electro-production of vector mesons

I.V. Anikin

Bogoliubov Laboratory of Theoretical Physics, JINR, Russia

O.V. Teryaev

Bogoliubov Laboratory of Theoretical Physics, JINR, Russia

## Abstract

We present an analysis of genuine twist-$`3`$ quark contributions to the amplitude of exclusive electro-production of transversely polarized vector mesons. Using the formalism based on the momentum representation we calculated all the genuine twist-$`3`$ terms of quark contributions to $`\gamma _T^{}\rho _T`$ amplitude. We found that these terms can not be factorized owing to the existence of the infrared divergencies in the amplitude of hard sub-processes.

The electro-production of transverse polarized vector mesons $`\mathrm{hadron}(p_1)+\gamma ^{}(q)\rho (p)+\mathrm{hadron}(p_2)`$ provides the well-known example of QCD factorization breaking. Indeed, the factorization is valid in the case of longitudinally polarized $`\rho `$-meson production unless the high twist effects are included . However, the description of transversely polarized meson production is strongly complicated due to the existence of infrared divergencies in amplitudes, breaking down the factorization (see e.g. ,  and references therein).

The amplitude of transverse vector meson production corresponds to the $`1/Q`$ suppressed contributions in comparison with the longitudinal vector meson case . At the same time, the resent experiments show that the transverse mesons production amplitudes do sensible contributions even at moderate virtualities $`Q^2`$ . So, to describe these processes we are forced to deal with the taking into account the terms of $`1/Q`$ – power while, on the other hand, it should yield to the problems related to the factorization theorem. In Ref. , the analysis of twist-$`2`$ amplitudes for hard exclusive lepto-production of mesons in terms of generalized parton distributions (GPDs or SPDs) was presented. Recently authors of  discussed the factorization problems in the electro-production of light vector mesons from transversely polarized photons. They have taken into account just kinematical twist-$`3`$ terms. Also the helicity flip amplitude of transversely polarized vector mesons production was considered within the kinematical Wandzura–Wilczek approximation .

Although the kinematical and dynamical (genuine) higher twists contributions are, generally speaking, independent, there are notable exceptions in some kinematical regions. In Deep Inelastic scattering at $`x_B1`$ the kinematical higher twist terms, described by Nachtmann variable, lead to inconsistencies unless the genuine higher twists are taken into account . One cannot exclude, that the genuine higher twists may cure, at least partially, the problem arising from the treatment of the end-point regions in hard electro-production.

Thus we want to study the role of genuine twist-$`3`$ contributions in the factorization theorem breaking. We adhere the approach based on the momentum representation the basic stages of which are expounded in previous papers -.

We now derive the electro-production amplitude of transversely polarized $`\rho `$-meson in terms of the coefficient function and ”soft” functions. First let us introduce the kinematics in such a manner: the $`p`$ is momentum of transversely polarized $`\rho `$-meson and its polarization vector is $`e^T`$; the momentum of virtual photon is denoted by $`q(Q^2=q^2)`$. We assume that the initial hadron momentum $`p_1`$ and final hadron momentum $`p_2`$ are collinear, and $`p_1^2=p_2^2=t=0`$, neglecting all the relevant higher twists contributions. In addition we neglect squares of meson masses, restricting ourselves to twist 3 contributions. With the help of $`p_1`$ and $`p_2`$ momenta we build the general relative momentum $`\overline{P}=(p_2+p_1)/2`$ and transfer momentum $`\mathrm{\Delta }=p_2p_1`$. In our approximation $`\overline{P}^2=\mathrm{\Delta }^2=0`$.

We introduce the parameterizations of all the matrix elements needed for calculations of amplitudes. We use the axial gauge $`nA=0`$, where $`n^2=0,np=1`$ In terms of the light-cone basis vectors the $`\rho `$-meson–to–vacuum matrix elements can be written as (keeping just the terms up to the twist-$`3`$ order):

$`p|\overline{\psi }(0)\gamma _\mu \psi (z)|0\stackrel{}{=}\phi _1(y)p_\mu +\phi _3(y)e_\mu ^T,p|\overline{\psi }(0)\gamma _\mu \stackrel{}{_\rho ^T}\psi (z)|0\stackrel{}{=}\phi _1^T(y)p_\mu e_\rho ^T,`$ (1)
$`p|\overline{\psi }(0)\gamma _5\gamma _\mu \psi (z)|0\stackrel{}{=}i\phi _A(y)\epsilon _{\rho \alpha \beta \delta }e_\alpha ^Tp_\beta n_\delta ,`$ (2)
$`p|\overline{\psi }(0)\gamma _5\gamma _\mu \stackrel{}{_\rho ^T}\psi (z)|0\stackrel{}{=}i\phi _A^T(y)p_\mu \epsilon _{\rho \alpha \beta \delta }e_\alpha ^Tp_\beta n_\delta ,`$ (3)
$`p|\overline{\psi }(0)\gamma _\mu gA_\rho ^T(z_2)\psi (z_1)|0\stackrel{}{=}\mathrm{\Phi }(y_1,y_2)p_\mu e_\rho ^T,`$ (4)
$`p|\overline{\psi }(0)\gamma _5\gamma _\mu gA_\rho ^T(z_2)\psi (z_1)|0\stackrel{}{=}iJ(y_1,y_2)p_\mu \epsilon _{\rho \alpha \beta \delta }e_\alpha ^Tp_\beta n_\delta ,`$ (5)
where $`\stackrel{}{=}`$ denotes the Fourier transformation with measure ($`z_i=\lambda _in`$)

$`dye^{iypz}`$ $`\mathrm{for}\mathrm{quark}\mathrm{correlators},`$ (6)
$`dy_1dy_2e^{iy_1pz_1i(1y_2y_1)pz_2}`$ $`\mathrm{for}\mathrm{quark}\mathrm{gluon}\mathrm{correlators}.`$ (7)
In the adopted collinear limit we directly obtain that $`p_2=(1\xi )/(1+\xi )p_1=\kappa p_1`$. Therefore, $`u(p_2)=\sqrt{\kappa }u(p_1)`$. This lead to the symmetrical nucleon spinor forms like $`\overline{u}(p_1)\sqrt{\kappa }\widehat{n}u(p_1)`$ whose calculation is straightforward. As a result we are able to introduce the parameterization of relevant matrix elements for nucleons. Further, keeping the twist-$`2`$ terms only, we write

$`p_2|\overline{\psi }(0)\gamma _\mu \psi (\stackrel{~}{z})|p_1\stackrel{}{=}H_1(x)\overline{P}_\mu `$ (8)
for the nucleon–nucleon matrix element of pure quark correlator .

In this paragraph we compute the quark contributions to the production amplitude which are generated by the diagrams with the quark legs coming from the nucleon blob. The calculating of the simplest Feynman diagrams give the following expression:

$`𝒜_{1,\mu }^{(q),\gamma _T^{}\rho _T}`$ $`=`$ $`4{\displaystyle \frac{C_F}{N_c}}{\displaystyle \frac{e_\mu ^T}{Q^2}}{\displaystyle _1^1}𝑑xH_1(x)\left[{\displaystyle \frac{1}{x\xi +iϵ}}{\displaystyle \frac{1}{x+\xi iϵ}}\right]`$ (10)
$`{\displaystyle _0^1}{\displaystyle \frac{dy}{y(1y)}}\left(\phi _A(y)+\phi _3(y)\right).`$
Next, summing the diagrams that do contribute to the genuine twist-3, we finally obtain the expression for the quark distribution amplitude. We write

$`𝒜_{2,\mu }^{(q),\gamma _T^{}\rho _T}`$ $`=`$ $`4{\displaystyle \frac{C_F}{N_c^21}}{\displaystyle \frac{e_\mu ^T}{Q^2}}\{\xi {\displaystyle _1^1}dxH_1(x)[{\displaystyle \frac{1}{(x+\xi iϵ)^2}}+{\displaystyle \frac{1}{(x\xi +iϵ)^2}}]_1^{(q)}+.`$ (11)
. $`{\displaystyle _1^1}dxH_1(x)[{\displaystyle \frac{1}{x\xi +iϵ}}{\displaystyle \frac{1}{x+\xi iϵ}}]_2^{(q)}\},`$ (12)
where

$`_1^{(q)}={\displaystyle _0^1}dy_1dy_2\left\{\stackrel{~}{J}(y_1,y_2)\right[4C_F({\displaystyle \frac{1}{(1y_1)^2}}{\displaystyle \frac{1}{(1y_2)^2}})+..`$ (13)
$`..`$ $`C_A({\displaystyle \frac{1}{y_2(1y_1)}}{\displaystyle \frac{1}{y_1(1y_2)}})+{\displaystyle \frac{2C_FC_A}{y_1+y_2}}({\displaystyle \frac{1}{y_1}}{\displaystyle \frac{1}{y_2}})].`$ (14)
. $`\stackrel{~}{\mathrm{\Phi }}(y_1,y_2)\left[4C_F\right({\displaystyle \frac{1}{(1y_1)^2}}+{\displaystyle \frac{1}{(1y_2)^2}})+C_A({\displaystyle \frac{1}{y_2(1y_1)}}+{\displaystyle \frac{1}{y_1(1y_2)}})+..`$ (15)
$`..`$ $`{\displaystyle \frac{2C_FC_A}{y_1+y_2}}({\displaystyle \frac{1}{y_1}}+{\displaystyle \frac{1}{y_2}})]\};`$ (16)
$`_2^{(q)}={\displaystyle _0^1}dy_1dy_2\left\{\stackrel{~}{J}(y_1,y_2)\right[C_F({\displaystyle \frac{1}{y_2(1y_1)}}{\displaystyle \frac{1}{y_1(1y_2)}})+{\displaystyle \frac{2C_FC_A}{y_1+y_2}}({\displaystyle \frac{1}{y_2}}{\displaystyle \frac{1}{y_1}})].`$ (17)
. $`\stackrel{~}{\mathrm{\Phi }}(y_1,y_2)[C_F({\displaystyle \frac{1}{y_2(1y_1)}}+{\displaystyle \frac{1}{y_1(1y_2)}})+{\displaystyle \frac{2C_FC_A}{y_1+y_2}}({\displaystyle \frac{1}{y_2}}+{\displaystyle \frac{1}{y_1}})]\}.`$ (18)
Here we introduced the notation for the general parameterizing functions $`\stackrel{~}{\mathrm{\Phi }}`$ and $`\stackrel{~}{J}`$ consisting of the kinematical and genuine (dynamical) twist-$`3`$ parts, i.e., for example,

$`\stackrel{~}{\mathrm{\Phi }}(y_1,y_2)=\phi _1^T(y_1)\delta (y_1y_2)+\mathrm{\Phi }(y_1,y_2).`$ (19)
The structure integrals of the quark distribution amplitude possess poles of second order also. This fact lead to the vulnerability of factorization theorem unless the functions $`\stackrel{~}{\mathrm{\Phi }}`$ and $`\stackrel{~}{J}`$ vanish at $`y_i1`$ or $`y_i0`$ more rapidly then the first power of $`(1y_i)`$ or $`y_i`$. Within our formalism we have reproduced, also, the results for the kinematical twist 3 gluon distribution amplitude obtained in .

In summary, we have computed the total quark contributions to the transversely polarized $`\rho `$-meson electro-production. We have found the novel singular terms of the genuine twist-$`3`$ to the quark contribution amplitude, which must be included in the general analysis of problems related with the factorization theorem breaking.

Also, we would like to thank to A.V. Efremov, N. Kivel, A. Radyushkin, A. Schaefer, C. Weiss for fruitful discussions and comments. This work has been supported in part by INTAS (Project 587, call 2000) and by RFFI Grant 00-02-16696.